Classical Waves

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Classical Waves

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David J. Pine

and

T. C. Lubensky

c Copyright 2014

D. J. Pine & T. C. Lubensky

December 24, 2014

Contents

List of tables x

Part I Oscillations 1

1 Oscillations 3

1.1 The simple pendulum 3

1.2 Diatomic molecule 7

1.3 The ideal spring 10

1.3.1 The horizontal spring 11

1.3.2 The vertical spring: an inhomogeneous ODE 11

1.3.3 Initial conditions and constraints 14

1.4 Energy of a simple harmonic oscillator 15

1.5 Complex exponential notation 19

1.5.1 Complex exponentials 19

1.5.2 Complex exponentials in harmonic oscillator problems 21

1.6 Series expansions 23

1.7 Summary of important points of Chapter 1 25

Problems 27

2 Damped Oscillations 31

2.1 The damped pendulum 31

2.1.1 Underdamped response: ( /2)2 < !20 34

2.1.2 Overdamped response: ( /2)2 > !20 37

2.1.3 Critically damped oscillations: ( /2)2 = !20 38

2.1.4 Summary of damped oscillations 39

2.2 Energy loss from damping 40

2.3 Shock absorbers 43

Problems 48

3 Resonance 52

3.1 Forced oscillators and resonance 52

3.1.1 Resonance 53

3.1.2 Response functions 59

3.1.3 Dissipation 61

iii

t

iv Contents

3.1.5 Transients 63

3.1.6 Indirect forcing 65

3.2 Resonant filters and detectors 68

3.2.1 Radio receiver 70

Problems 74

4 Normal Modes 78

4.1 Systems with two degrees of freedom 78

4.1.1 Two coupled pendulums 78

4.1.2 Normal modes of two coupled pendulums 82

4.1.3 Weak coupling 86

4.1.4 Energy in normal modes 88

4.2 Matrix formulation of normal modes 89

4.2.1 Two equal-mass particles 90

4.2.2 The meaning of normal coordinates 93

4.2.3 Specific solutions 96

4.2.4 Three masses on a string 98

4.2.5 Systems of unequal masses 102

4.2.6 Geometry and symmetry 107

4.3 Normal modes of carbon dioxide 108

4.3.1 Longitudinal modes 108

4.3.2 Transverse modes 112

4.4 Damping and normal modes 115

4.5 Forced oscillations 118

4.5.1 Steady-state response 118

4.6 Summary of important points of Chapter 4 121

Problems 122

5.1 Identical point masses on a massless string 128

5.1.1 Extrapolation from small to large N solutions 130

5.1.2 Symmetry and the matrix formulation 137

5.1.3 Matrix diagonalization and similarity transformations 138

5.1.4 Numerical schemes for matrix diagonalization 139

5.1.5 Jacobi method of matrix diagonalization 139

5.2 Sequences of dierent point masses on a string 144

5.2.1 Periodic alternating sequence of masses 144

5.2.2 Random sequences of masses: Localization 153

5.2.3 A Fibonacci sequence of two masses: Quasiperiodicity 159

5.2.4 Structure: The spatial distribution of masses 161

5.3 Traveling waves 162

5.4 Pulse propagation and dispersion 165

5.5 The continuum approximation 165

t

v Contents

Problems 173

6 Strings 175

6.1 Traveling waves 175

6.2 Normal modes 175

6.3 Interfaces 175

6.4 Energy transport 175

6.5 Infinite strings 175

Problems 176

7.1 Fourier series 178

7.1.1 Sine and cosine Fourier series 178

7.1.2 Fourier series in time 183

7.1.3 Basis functions and orthogonality 183

7.1.4 Power in periodic waves 184

7.1.5 Complex exponential Fourier series 186

7.2 Fourier transforms 187

7.2.1 Fourier transform pairs 187

7.2.2 Some examples 190

7.2.3 The uncertainty principle 192

7.2.4 Some theorems 194

7.3 Discrete Fourier Transforms 195

Problems 201

8.1 Sound 206

8.2 Reflection and refraction 206

8.3 Electromagnetic waves 206

Problems 207

9.1 Interference 208

9.2 Diraction 208

9.3 Fourier optics 208

9.4 Scattering 208

Problems 209

t

vi Contents

A.1 Vectors 213

A.2 Matrices 214

A.2.1 Matrix operations 214

A.3 Properties of matrices 215

A.3.1 Determinant of a matrix 215

A.3.2 Inverse of a matrix 216

A.3.3 Matrix symmetry 216

A.3.4 Hermitian matrices 217

A.3.5 Positive definite matrices 217

A.3.6 Orthogonal matrices 217

A.3.7 Matrix identities 217

A.4 Transformations 218

A.4.1 Rotations 218

A.4.2 Choleski decomposition 218

A.4.3 Similarity transformations 220

A.5 Eigenvalue equations 220

Notes 233

Illustrations

1.2 Sinusoidal trajectory 5

1.3 Approximations for sin . 6

1.4 Potential energy for a diatomic molecule 8

1.5 Ideal spring on a frictionless surface 11

1.6 Vertical spring 12

1.7 Mass with spring on inclined plane 13

1.8 Colliding masses with springs 15

1.9 Energies for mass on a spring as a function of position 17

1.10 Stretched springs 18

1.11 The complex plane 20

1.12 Oscillator on Inclined Plane 29

1.13 Four Springs 30

1.14 Drooping sprints 30

2.1 Damped pendulum 32

2.2 Transient response of an underdamped oscillator 36

2.3 Transient response of an overdamped oscillator 37

2.4 Transient response for damped oscillators 40

2.5 Automobile suspension system with shock absorber 43

2.6 Response of automobile shock absorbers 46

2.7 Torsion Oscillator 50

3.1 Damped forced mass suspended from a spring 53

3.2 Amplitude and phase of sinusoidally driven oscillator 56

3.3 Trajectories of oscillators driven at and o resonance 58

3.4 Response functions of an oscillator 59

3.5 Resonant power absorption of an oscillator 61

3.6 Transient response of a driven oscillator 66

3.7 Pendulum with a moving pivot point 67

3.8 Amplitude and phase of shaking pendulum 69

3.9 Amplitude vs. frequency for seismometer pendulum 70

3.10 Filtered and unfiltered radio signal 71

3.11 Resonant circuit with a resistor, inductor, and capacitor in series. 72

3.12 Resonant RLC circuit. 76

4.1 Two identical pendulums coupled by a spring 79

4.2 Modes of two pendulums coupled by a spring. 83

4.3 Displacements of two strongly coupled pendulums 86

vii

t

viii Illustrations

4.5 Horizontal masses on a frictionless surface coupled by two springs 90

4.6 Normal coordinates for 2-mass 2-spring problem. 95

4.7 Three masses on a string under tension 98

4.8 Normal modes of 3 masses on a string 101

4.9 Trajectories of weakly-coupled pendulums of unequal masses 106

4.10 Normal mode coordinate transformation 107

4.11 Longitudinal and transverse vibrations of CO2 109

4.12 Normal modes of CO2 111

4.13 Transverse modes of CO2 : linear and circular polarization 114

4.14 Resonance curves for two weakly coupled pendulums with forcing 121

4.15 Wilberforce pendulum 125

5.1 N masses on a string 129

5.2 Normal modes for N masses on a string for N = 1, 2, 3. 131

5.3 Normal frequencies vs. mode number for identical masses on a string 133

5.4 Normal modes for five particles on a string 135

5.5 Normal modes for 32 particles on a string 145

5.6 Normal modes for alternating sequence of 32 masses 146

5.7 Normal modes 1, 2, 16, and 17 of 32 alternating masses on a string 147

5.8 Normal modes 30-32 for 32 alternating masses on a string 148

5.9 Coordinates for alternating masses on a string 148

5.10 Frequencies of normal modes for alternating masses on a string 151

5.11 Normal modes for a random sequence of particles 154

5.12 Normal modes for a dierent random sequence of particles 155

5.13 High frequency normal modes for a random sequence of particles 156

5.14 Localization length vs. frequency 157

5.15 Normal modes for a random sequence of 32 spheres 158

5.16 Density of states for a random sequence of masses 159

5.17 Frequencies of normal modes for a Fibonacci sequence 160

5.18 Normal modes for a Fibonacci sequence of two masses 161

5.19 Normal modes for an 89-mass Fibonacci sequence of two masses 162

5.20 Structure factors of dierent sequences of masses 163

5.21 Discrete approximation for second derivative. 166

5.22 Gaussian wave packet propagating in a nondispersive medium 167

5.23 Superposition of waves of with dierent velocities 169

5.24 Propagation of Gaussian wave packet with dispersion 171

5.25 Gaussian wave packet at three times 172

7.1 Periodic sawtooth and triangle waveforms 178

7.2 Cosine and sine waves for the first three terms in a Fourier series 179

7.3 Fourier series of sawtooth function 181

7.4 Fourier components of sawtooth function 182

7.5 Sinusoidal wave packet with Gaussian envelope 188

7.6 Fourier transform of sinusoidal wave packet with Gaussian envelope 190

7.7 (a) Square pulse and (b) its Fourier transform. 192

t

ix Illustrations

7.9 Discrete Fourier transform of an exponential wave packet 199

7.10 Filtering a noisy Gaussian pulse 201

B.1 Double pendulum. 227

Tables

7.1 Ordering of spatial frequencies for the discrete Fourier transform. 198

x

PART I

OSCILLATIONS

1 Oscillations

Everything oscillates: stars revolving in galaxies, planets orbiting stars, seas, clocks, hearts

the list is endless. The space all around you is filled with radio waves that oscillate. There

are theories that propose that the universe itself oscillates. Even when you think youve

found something that doesnt oscillatesay the chair you are sitting onyou find that it

is made of molecules and atoms that oscillate, and the atoms themselves are made up other

things that oscillate: electrons, protons, and neutrons. When we look even deeper into what

makes up these tiny constituents, we find more things that oscillate: photons, quarks, and

maybe even strings. So if our goal is to understand the world around us, the world within

us, and the worlds beyond us, we need to understand oscillations.

In this chapter, we develop the kinematics and dynamics of sinusoidal oscillations. We

do so by first considering two physical examples: an oscillating simple pendulum and a

vibrating diatomic molecule. We then turn to a common example, a mass attached to an

ideal spring. We also begin what will be an ongoing task of developing mathematical tools

for describing oscillations. In this chapter, we introduce complex exponentials, a wonderful

mathematical tool that greatly simplifies the algebra of sinusoidal functions.

We begin our discussion of oscillations with the simple pendulum, a mass hung from a fixed

point so that it can freely swing back and forth. When pulled from its equilibrium position

and released, gravity causes the pendulum to fall and the pendulum begins oscillating. In

reality, the oscillations die out in time because of friction, but in this chapter we consider

only the idealized case where there is no friction. In Chapter 2, we introduce friction.

Our first task is to describe the motion of the pendulum quantitatively. To this end, we

reduce the pendulum to its simplest form, a mass m attached to one end of a rigid massless

rod of length l that pivots about its other end, as shown in Fig. 1.1. When the pendulum is

displaced from = 0, the rod is constrained to move along an arc described by s = l. The

component of gravity directed along the arc s provides a restoring force Fr back towards

= 0. It is given by

Fr = mg sin . (1.1)

To obtain the equation of motion for the pendulum, we apply Newtons second law along

3

t

4 Oscillations

t

Fig. 1.1 Simple pendulum, captured at an instant when it is moving towards = 0.

d2 s

m = mg sin . (1.2)

dt2

Using s = l to eliminate s, Eq. (1.2) yields a dierential equation for the angular displace-

ment of the pendulum

d2 g

+ sin = 0 . (1.3)

dt2 l

Equation (1.3) is not easily solved for arbitrary values of . However, in the limit of small

oscillations, that is, when 1, we can expand sin in a Taylor series,

1 3 1 5

sin = 3! + 5! ... , (1.4)

and obtain an analytic solution to Eq. (1.3) by keeping only the first term. Using the ap-

proximation sin ' , Eq. (1.3) becomes

d2

+ !20 = 0 , (1.5)

dt2

p

where we have defined the frequency !0 g/l. You can readily verify that the units for

!0 are [time] 1 , as g has units of acceleration [length]/[time2 ] and l has units of [length].

We pause briefly to make some general mathematical observations about Eq. (1.5). We

have written Eq. (1.5) so that all the terms involving the dependent variable appear on

the left hand side; any terms not involving should appear on the right hand side. This is

a standard form for writing dierential equations. Equation (1.5) is a homogeneous linear

second order ordinary dierential equation (ODE). It is homogeneous because the right

hand side is zero; there is no term in the equation that does not involve the dependent

variable . It is second order because the highest derivative appearing in the equation is a

second derivative. It is ordinary because it involves only one independent variable, t in this

case. It is linear because the dependent variable appears only as or as derivatives of ;

there are no terms like where , 1 or terms like (n) where , 0 and (n) denotes

t

5 The simple pendulum

t

Fig. 1.2 Sinusoidal trajectory of simple harmonic motion described by Eq. (1.6). The initial

displacement is a and the initial (angular) velocity is !0 b. The period is T = 2/!0 .

the nth derivative (n 1) of . We will make use of these mathematical properties, noting

when we do so, in this and the following chapters.

You probably know from your previous studies that the solutions to Eq. (1.5) are si-

nusoidal. In fact both sine and cosine functions are solutions. That there are two linear

independent solutions, sine and cosine, follows from the fact that Eq. (1.5) is a second

order linear ODE; third order linear ODEs have three linearly independent solutions, ...

Because Eq. (1.5) is linear, the sum of of any set of solutions is also a solution. This is the

principle of linear superposition. In particular, any sum of the sine and cosine solutions is

also a solution. Thus, the most general solution is

The solutions satisfy Eq. (1.5) only when the coefficient of t in the sine and cosine functions

is !0 . You can check this by substituting Eq. (1.6) into Eq. (1.5). The cosine and sine

functions are periodic with period 2, which means that the solution (t) is periodic with

period T = 2/!0 . Mathematically, we can express the periodicity of (t) by writing (t) =

(t + T ). The angular frequency !0 has units of radians per second and should not be

confused with the cyclic frequency 0 = 1/T measured in cycles per second. The two

frequencies are related to each other by !0 = 20 . As 2 has no units, both !0 and 0

have the same units of inverse time; radians and cycles are unitless.

The constants a and b are determined by the initial conditions. Setting t = 0, we obtain

the following equations for the initial position and velocity from Eq. (1.6)

0 = a !0 sin(!0 0) + b !0 cos(!0 0) = !0 b ,

(0) (1.8)

0

(t) = 0 cos !0 t + sin !0 t. (1.9)

!0

t

6 Oscillations

t

Fig. 1.3 Approximations for sin .

This equation can be expressed in a dierent but equivalent form that is often useful:

= A(cos !0 t cos + sin !0 t sin ) . (1.11)

Equating the coefficients of cos !0 t and sin !0 t in Eq. (1.9) and Eq. (1.10) leads to

q

0

A = 02 + (0 /!0 )2 ; = tan 1 . (1.12)

!0 0

Figure 1.2 shows a typical function (t) in which both 0 and 0 , which is the slope of (t)

at t = 0, are positive. It is clear from Eq. (1.10) that (t) reaches its maximum value of A at

time t = /!0 = T/2.

The motion of an oscillator whose equation of motion has the same form as Eq. (1.5)

is often called simple harmonic motion because its trajectory is described by a simple

sinusoidal (harmonic) function. From the point of view of mathematics, it is certainly the

simplest kind of oscillator to treat. It is important to keep in mind, however, that to obtain

this simple result, we had to replace the true restoring force given by Eq. (1.1) with the

linearized approximate force

Fr` = mg . (1.13)

Therefore, our treatment is an approximation that is valid only when sin ' , that is, when

1. We could obtain more accurate results by including more terms from Eq. (1.4) in

our approximation of the force law:

Fr ' mg 3!1 3 + 5!1 5 ... (1.14)

2 4

= mg 1 3!1 + 5!1 ... . (1.15)

When a force law contains terms that are nonlinear in , the force is said to be anharmonic.

In this case, the lowest order anharmonic term is cubic. The anharmonic terms can be ig-

nored to the extent that they are small compared to the linear, or harmonic, term. Examining

Fig. 1.3, it appears that the linear approximation is good out to about 0.5 rad. Including

the first anharmonic term gives an adequate approximation out to about 1 rad. If we

t

7 Diatomic molecule

limit the amplitudes of oscillation to be less than 0.5 rad or about 30 , the relative error in

the force incurred by keeping only the linear term is approximately

h i

F Fr Fr` mg 1 3!1 2 + 5!1 4 ... 1

= ' 16 2 ' 0.04 , (1.16)

Fr Fr` mg

or about 4%. Equation (1.16) also implies that reducing the amplitude of the pendulum by

another factor of two reduces the error in the force by a factor of four. So the harmonic

approximation can be made as accurate as we please provided we limit our description to

sufficiently small amplitudes of the pendulum.

In the analysis of oscillating systems, one rarely encounters a purely linear force law

for displacements from the equilibrium position. However, expanding the force law in a

Taylor series and keeping only the first term, as we did above for the case of the pendulum,

is almost always the first line of attack in trying to describe an oscillating system. Doing

so allows us to solve the equation of motion analytically. Including higher order terms in

the Taylor series expansion of sin requires analytical approximation methods or numer-

ical techniques. The range of displacements over which such an approximation is valid

depends on the force law that describes the physical system and the accuracy required of

the solution. In any real system, you should always ask yourself how much error is incurred

by keeping only the linear term by performing an analysis similar to that done above. In

the next section, we consider molecular vibrations, another system for which the harmonic

approximation is often made.

The atoms in a diatomic molecule vibrate back and forth, alternately stretching and com-

pressing the bond between atoms, as the system exchanges kinetic and potential energy.

The interatomic potential energy is complicated but is sometimes approximated by the

Morse potential, a function given by

2

U(r) = U0 1 e (r r0 )/ , (1.17)

where r = R2 R1 is the distance between atoms, as illustrated in Fig. 1.4. The potential

energy well has a minimum at r = r0 , the equilibrium distance between the atoms, and has

a width of approximately . Figure 1.4 shows a plot of Eq. (1.17) for the case = 0.1r0 .

For r < r0 , which corresponds to the two atoms being closer than the equilibrium distance,

the potential is strongly repulsive due to the electron clouds around each atom. For r r0 ,

the potential increases from its minimum of zero at r = r0 towards a value of U0 . This is the

binding energy between the two atoms; the atoms become unbound if the kinetic energy

associated with the vibrational motion of the atoms exceeds this value.

When the separation between the two atoms in the molecule deviates from r = r0 , the

atoms oscillate back and forth. Analyzing the oscillations about the equilibrium position

at r = r0 leads to a nonlinear dierential equation that cannot be solved analytically if we

t

8 Oscillations

U(r)

U0

R

R1 R2

U0

R

R1 R2

t

0

0.9 1.0 1.1 1.2 1.3 r/r0

Fig. 1.4 Potential energy for a diatomic molecule based on Eq. (1.17) with = 0.1r0 . The solid

trace shows the potential energy U(r) between two covalently bound atoms as a

function of the distance r R2 R1 between them. The potential is a minimum at the

equilibrium separation r0 . The dashed line shows the harmonic approximation U s ( r)

where r r r0 . The coordinates of the atomic nuclei are shown on the right

together with a crude indication of the electron clouds. The two atoms, located at

R = R1 and R = R2 , are shown in their compressed and stretched configurations in the

upper and lower diagrams, respectively

use the full potential given by Eq. (1.17). To simplify the analysis, the potential can be

expanded in a Taylor series about r = r0 where the potential is a minimum (zero in this

case):

U 00 (r0 ) U (n) (r0 )

U(r) = U(r0 ) + U 0 (r0 )(r r0 ) + (r r0 )2 + ... + (r r0 )n + ... (1.18)

2! n!

Applying this formula to Eq. (1.17) while noting that U(r0 ) = 0 and U 0 (r0 ) = 0 gives

" 2 #

r r3 7 r4

U(r) = U0 2 3

+ + ... , (1.19)

12 4

where r r r0 is the displacement from the equilibrium distance between the atoms.

Keeping only the first term in the series expansion results in a potential energy that is

quadratic in the displacement r. In this case, the potential energy can be written in a form

equivalent to that of an ideal spring,

1

U s ( r) = k r2 . (1.20)

2

Comparing Eq. (1.20) to Eq. (1.19) we see that the spring constant is given by

2

k = 2U0 / . (1.21)

The dashed line in Fig. 1.4 shows U s ( r). It is apparent from the plot that the approximation

is accurate only if the oscillations are quite small, something like a few per cent of the

distance r0 between the atoms. You can verify that the units of k are those of a spring,

t

9 Diatomic molecule

[force]/[length], as U0 is an energy with units of [force][length] (i.e. work) and r has units

of [length].

The force law corresponding to the potential given by Eq. (1.21) can be obtained from

the gradient (spatial derivative) of the potential:

@U s (r)

Fs = = k r. (1.22)

@r

We recognize Eq. (1.22) as the usual linear spring force law where the force is proportional

to minus the displacement from the equilibrium position. This force law has the same form

as Eq. (1.13), the linearized force law for the pendulum. In both cases, the restoring force

(note the minus sign) is proportional to the displacement. For the pendulum the spring

constant" is mg/l, while for the oxygen molecule it is 2U0 / 2 . Because the force laws for

the pendulum and the diatomic molecule are the same in the limit of small displacements,

we expect their behavior to be the same, at least for small amplitudes.

To find explicit equations for the motion of the atoms, we need to write down Newtons

second law. Lets assume that the two atoms are the same element, say oxygen, so both

atoms have the same mass m. Referring to Fig. 1.4 for the coordinates, an equation of

motion can be written for each particle

d 2 R1

F1 = m = k(R2 R1 r0 ) (1.23)

dt2

d 2 R2

F2 = m 2 = k(R2 R1 r0 ) (1.24)

dt

We are looking for an equation that describes the vibration of the two atoms, that is, an

equation that governs how the distance R2 R1 between the two atoms oscillates in time.

To this end, we subtract Eq. (1.23) from Eq. (1.24), which gives

d2 (R2 R1 )

m = 2k(R2 R1 r0 ) . (1.25)

dt2

Instead of working with R2 R1 , we subtract o the equilibrium separation r0 between the

two atoms and define r = R2 R1 r0 , which at equilibrium is zero. With this definition

Eq. (1.25) becomes

d2 r

= k r, (1.26)

dt2

where we have made use of the fact that r0 is a constant so its derivative is zero. In obtaining

Eq. (1.26), we have reduced the two-body problem (with two masses and two equations)

to a one-body problem with an eective masspof . The constant m/2 and is called the

reduced mass of the system. Defining !0 = k/, Eq. (1.26) can be rewritten as

d2 r

+ !20 r = 0 , (1.27)

dt2

which has the same form as Eq. (1.5) and therefore has the same solution with a sinusoidal

form given

p by Eq. (1.6). This means that the diatomic atom oscillates with a frequency

!0 = k/, where k is given by Eq. (1.21).

t

10 Oscillations

Looking back over our work, we see that the spring constant k is the second derivative

of the potential evaluated at zero displacement:

@2 U

k= , (1.28)

@x2 x=0

where x is the (generalized) displacement from the equilibrium position. For the diatomic

molecule x = r; for the pendulum x = . This follows from the fact that the potential

energy U = 12 kx2 + terms proportional to higher powers of x. Viewed from this perspective,

k is simply the curvature of the potential at zero displacement.

Our analysis shows that approaching the problem of oscillations starting from conserva-

tion of energy yields the same results as does approaching the problem starting from forces

and Newtons second law. Evidently, keeping only terms up to the quadratic term in the

energy formulation is equivalent to keeping the linear term in the force formulation.

An important point to take away from the analysis is that physical systems that can be

described by the same dierential equations exhibit the same physical behavior. This is a

very powerful idea because it means that once you understand one system, for example

the pendulum as a harmonic oscillator, you understand the essential physics of many sys-

tems. Here we discover that diatomic molecules also behave like harmonic oscillators if

the displacement from the equilibrium position is sufficiently small. In fact, as we shall

see, nearly all vibrating systems behave like harmonic oscillators provided the amplitude

of the oscillations is sufficiently small.

Exercise 1.2.1 Over what range of distances is the quadratic approximation for Eq.

(1.17) valid to within 10% of the exact version? Use parameters for O2 : bond length

r0 = 0.121 nm, = 0.1r0 , and bond energy U0 = 5.2 eV. Hint: Using the Taylor

expansion of U(r) given by Eq. (1.19), find an equation for U(r)/U(r) [U(r)

U s (r)]/(U(r).

Nothing in life is perfect, so the saying goes, yet physicists are forever analyzing the ideal

spring. As in the rest of life, real springs do not achieve our ideal. In particular, they do not

obey a simple linear force law, except to some degree of approximation that is valid for suf-

ficiently small displacements. In this way, real springs are just like pendulums and diatomic

molecules. Nevertheless, the ideal spring serves a launching point for most discussions of

oscillations.

Our purpose here in not having started our discussion of oscillations with the ideal spring

is to emphasize that the harmonic approximation is just that, an approximation that is

generally valid only over some finite range of displacements from the equilibrium position.

With that fact firmly in mind, we briefly reformulate the problem of small oscillations in

terms of the ideal spring and recapitulate the main results of the previous sections.

t

11 The ideal spring

We adopt the notation usually associated with an ideal spring with the familiar linear force

law (Hookes Law)

F = kx , (1.29)

where x is the distance the spring is stretched (see Fig. 1.5). The spring is attached to a

mass m that moves horizontally on a frictionless surface, as shown in Fig. 1.5 (recall that

we are treating only frictionless cases in this chapter). Applying Newtons second law to

the system, we obtain the equation of motion

F = kx = m x , (1.30)

which can be rewritten in the familiar form

d2 x

+ !20 x = 0 , (1.31)

dt2

p

where !0 = k/m. Equation (1.31) has the same form and the same solutions as the

linearized equation of motion Eq. (1.5) for the pendulum. In particular, the general solution

to Eq. (1.31) can be written as

x(t) = a cos !0 t + b sin !0 t , (1.32)

or equivalently as

x(t) = d cos(!0 t ); , (1.33)

where a, b, d, and can be determined from the initial conditions.

Suppose a mass hangs vertically from a spring attached to the ceiling as shown in Fig. 1.6.

If the length of the unstretched spring is L0 and y = L L0 is the distance the spring is

stretched downward from this position, then the spring force ky on the mass is upward

and negative, while the gravity force mg is downward and positive. The equation of motion

is

d2 y

m = ky + mg . (1.34)

dt2

t

Fig. 1.5 Ideal spring on a frictionless surface for unstretched spring (top) and stretched sprint

(bottom).

t

12 Oscillations

t

Fig. 1.6 (a) Displacement of a mass hanging from a vertical spring for the unstretched (left)

and stretched positions (right). (b) Displacement z from the equilibrium stretched

position y0 . Both y and z are positive in the figure.

d2 y k

+ y=g. (1.35)

dt2 m

p

Setting !0 = k/m, the left hand side of this equation is the same equation as Eq. (1.31)

above, but now it is set equal to a constant term rather than zero on the right hand side.

Writing an ODE so that all occurrences of the dependent variable y are on the left hand

side and all other terms are on the right hand side, an ODE is said to be homogeneous if

the right hand side is zero and inhomogeneous if the right hand side is a constant or any

function of the independent variable, which in this case is t. Thus, Eq. (1.35) is seen to be

inhomogeneous.

The solution y(t) = a cos !0 t + b sin !0 t that worked for the homogeneous case does not

work here because of the inhomogeneous term. In general, solving inhomogeneous ODEs

is more difficult than solving homogeneous ODEs, but in this case a simple trick suffices to

solve the equation. To begin, we consider the case when the mass is simply at rest hanging

at its equilibrium position. In that case, there is no acceleration and Eq. (1.35) becomes

k

y=g. (1.36)

m

Solving for y, we obtain the equilibrium position y0 of the mass

mg

y0 = . (1.37)

k

We now define a new variable z = y y0 = y mg/k, which is the displacement of the

mass from its equilibrium position. This means that y = y0 + z = mg/k + z. Substituting

this expression for y into Eq. (1.35) yields

d2 z k

+ z=0, (1.38)

dt2 m

which is the same homogeneous equation of motion we obtained for the horizontal spring.

Therefore, the same general solution works: z(t) = a cos !0 t+b sin !0 t, which is equivalent

to y(t) = mg/k + a cos !0 t + b sin !0 t.

t

13 The ideal spring

The trick for obtaining a homogeneous equation is to define the displacement z as the

displacement from the equilibrium position of the system, not from the position corre-

sponding to the unstretched spring. This illustrates a general principle for analyzing os-

cillations: always define your dynamic dependent variable as the displacement from the

equilibrium position. Because the system oscillates about its equilibrium position, this will

always result in the simplest equation of motion. If you look back at our formulation of the

problem of oscillations of a diatomic molecule, you will see that we used this trick without

much fanfare to eliminate r0 from Eq. (1.26), which otherwise would have resulted in an

inhomogeneous equation.

Example 1.1 Two stacked blocks of mass m rest on a frictionless inclined plane as

shown in Fig. 1.7. The bottom mass is connected to the wall by a spring of unstretched

length l0 and spring constant k. Calculate the initial rest position x0 of the two stacked

masses relative to the wall.

At time t = 0 the top mass is suddenly removed. Calculate the position x(t) for times

t > 0 and determine the maximum value of x(t). The spring is stretched when x > l0 ; the

velocity is positive when the mass moves upward along the incline.

Solution

To find the initial rest position x0 of the two masses, we balance the forces in the x direction

along the inclined plane, noting that the spring force is directed down the plane for a

stretched spring when x l0 > 0:

2mg

F x = k(x l0 ) 2mg sin = 0 ) x0 = l0 sin .

k

Thus, at equilibrium, x0 < l0 ; the spring is compressed and exerts an upward force that

balances the downward gravitational force of the two masses.

Now at t = 0, the upper mass is removed. The equilibrium position is no longer x0 , and

the remaining block will move. Its equation of motion is

d2 x

m = k(x l0 ) mgsin .

dt2

t

Fig. 1.7 (a) A spring attached to two masses on an inclined plane showing l0 the equilibrium

length of the spring and x0 it equilibrium position with the two masses. (b) The same

spring after one mass has been removed.

t

14 Oscillations

The initial conditions are x(0) = x0 and 3(0) = 0. We now follow the same procedure

outlined above for a vertical spring. First, calculate the new equilibrium position x1 by

setting the right hand side of the equation equal to zero: x1 = l0 (mg/k) sin . Next, define

z(t) = x(t) x1 . The equation of motion for z is now identical to Eq. (1.38) for a simple

harmonic oscillator with solution z(t) = a cos !0 t + b sin !0 t. The velocity at t = 0 is zero,

so b = 0, and z(0) = a = x0 x1 = (mg/k) sin . Putting this together gives

Most of the time, we use initial conditions to determine the two unknown constants in the

solution to the equation for a simple harmonic oscillator. However, these constants can be

determined in other ways. For example, position and velocity might be specified at some

time other than t = 0, or two positions might be specified at dierent times. Mathemati-

cally, we simply need two independent conditions about a systems motion. Once they are

set, it is only a matter of algebra to determine the unknown constants in the general solution

and thus determine the motion at all times.

Usually the physical conditions of a system, often in the form of constraints that are

released at t = 0, determine the initial conditions. For example, in the vertical spring

example of 1.3.2, the physical condition was that the mass was held in place (constrained

to be) at the rest length of the spring for times t < 0 and then released at time t = 0 so that

y(0) = 0 and y(0) = 0, implying z(0) = mg/k and z(0) = 0. In the example of the inclined-

spring, it is the compression of the spring by the two masses and the sudden removal of the

second mass that sets the constraint that x(0) = x0 and 3(0) = 0.

Collisions between a moving particle and a mass attached to a spring can lead to changes

in position and velocity and can thus set a system into oscillation. For example, consider a

particle of mass m moving with velocity 3 on a frictionless table; the mass m is moving to

the left initially so 3 < 0. At t = 0, it collides with a mass M at rest at the end of a spring

with force constant k, as shown in Fig. 1.8(a).

Consider first the case where the two particles stick together after they collide, as shown

in Fig. 1.8(b). The collision is completely inelastic, so energy is not conserved, but momen-

tum is. The velocity V 0 of the two particles at t = 0+ just after they collide is determined

by (m + M)V 0 = m3, which gives V 0 = m3/(m + M). The oscillator now has mass (m + M)

and natural frequency !1 = [k/(m + M)]1/2 . The motion of the mass for t > 0 is

m3

x(t) = sin !1 t , (1.39)

(m + M) !1

where x = 0 is the initial and equilibrium position of mass M.

In an elastic collision, both energy and momentum are conserved. The two masses are

in contact just for an instant. If the initial velocities of masses M and m are, respectively,

t

15 Energy of a simple harmonic oscillator

(a) (b)

k

M m M m

(c) (d)

V

t

k

M m M m

Fig. 1.8 (a) Mass m moving toward mass M before an inelastic collisions and (b) masses m

and M stuck together an moving with velocity after collision. (c) Masses M and m

approaching each other with respective velocities v and V before collision and (d)

masses m and M after the elastic collision.

m3 + MV = m30 + MV 0

2

(1.40)

1

2 m3 + 12 MV 2 = 12 m302 + 12 MV 02 ,

which leads to

M m 2m

V0 = V+ 3

M+m M+m (1.41)

2M m M

30 = V+ 3.

M+m m+M

If V = 0, the initial condition is V 0 = [2m/(M + m)] 3. In this case the oscillator has a

mass of M, since the two masses separate after the collision, and the natural frequency is

p

!2 = k/M.

As a final example, initial velocity may be determined by an impulse such as might be

delivered by a croquet mallet. An impulse is a force delivered over a short time period t:

Z t+ t

I= Fdt0 = p (1.42)

t

where p is the change in momentum brought about by the impulse. Thus the change in

velocity of a particle of mass m produced by an impulse is 3 = I/m.

Our analysis of the pendulum started by analyzing the force while our analysis of the

oxygen molecule started from the potential energy between the two atoms and from that

t

16 Oscillations

derived the force law. The two descriptions of oscillating systems are equivalent of course.

Here we take a closer look at simple harmonic oscillators from the perspective of the con-

servation of energy.

The change in potential energy U associated with compressing or stretching a spring

is defined as the negative of the work W done by the spring in moving from some position

x1 to another position x2 :

Z x2

1

U = W12 = kx dx = k(x22 x12 ) . (1.43)

x1 2

Taking the zero of potential energy to correspond to the unstretched spring (x = 0), the

potential energy of a spring becomes

1 2

U(x) = kx . (1.44)

2

The principle of conservation of energy for the mass-spring system can be obtained by

applying Newtons second law:

F = kx = m x , (1.45)

d3

kx = m . (1.46)

dt

Multiplying by dx and noting that dx = 3 dt, this can be written as

kx dx = m3 d3 (1.47)

1 2 1 2

m3 + kx = E , (1.48)

2 2

where E is an integration constant. The two terms on the left hand side in Eq. (1.48) are

the kinetic and potential energies, respectively. Equation (1.48) states that their sum is a

constant E, which we identify as the total energy of the system.

We can also verify explicitly that energy is conserved for this system. The time-dependent

displacement x(t) of the mass on the spring is sinusoidal, just as for the pendulum and di-

atomic molecule,

x(t) = a cos(!0 t ), (1.49)

p

where !0 = k/m. Writing 3 = x and noting that k = m!20 , we calculate the total energy

at any given time t to be

1 2 1 2

E= m3 + kx (1.50)

2 2

1 1

= m!0 a sin2 (!0 t

2 2

) + ka2 cos2 (!0 t ) (1.51)

2 2

1 h i

= m!20 a2 sin2 (!0 t ) + cos2 (!0 t ) (1.52)

2

1

= m!20 a2 . (1.53)

2

t

17 Energy of a simple harmonic oscillator

K(x)

U(x)

t

x

A A

Fig. 1.9 Kinetic energy K(x), potential energy U(x) and total energy E(x) = E as a function of

particle position.

Thus, the total energy of the system is explicitly demonstrated to be independent of time.

Incidentally we remark that the total energy of any linear oscillator can be written in the

form of Eq. (1.53) as a mass times the oscillation frequency squared times the amplitude

squared.

Of course, this means that at each point along its trajectory, the total energy of the os-

cillator is constant. Figure 1.9 shows plots of the kinetic, potential, and total energy as a

function of the oscillator position x. The kinetic energy is a maximum at the equilibrium

position x = 0 while the potential energy is a maximum at the turning points x = a,

where the pendulum momentarily comes to rest and reverses direction. The plot illustrates

that the sum of the potential and kinetic energies is constant, consistent with Eq. (1.48).

Exercise 1.4.1 pShow that the absolute value of the position-dependent velocity is given

by |3| = !0 a2 x2 and that the kinetic energy K(x) of the mass-spring system is

given by K(x) = 12 m!20 (a2 x2 ).

The potential energy U(x) of Eq. (1.44) is, of course, only the potential energy of the

stretched or compressed spring. The total potential energy will in general contain other

contributions. For example the vertical spring considered in 1.3.2 has an additional gravi-

tational potential energy Ug (y) = mgy (with a negative sign because y is measured down-

ward) to yield a total potential energy of

1 2

UT (y) = ky mgy, (1.54)

2

The potential energy here is at its minimum when dUT (y)/dy = 0, i.e.when y = y0 = mg/k.

An equilibrium state is one in which the force on the particle is zero or equivalently one

that is a minimum of the total potential energy.

Example 1.2 Figure 1.10 shows a block of mass m on a frictionless horizontal plane

attached to two identical springs with spring constants k and rest lengths l0 . The springs

are attached to fixed walls whose separation is such that in equilibrium, both springs are

stretched to a length 2l0 . What is the period of oscillation of the block: (a) if it is displaced

a distance x parallel to the springs and (b) if it is displaced a distance y perpendicular to

the plates?

t

18 Oscillations

t

Fig. 1.10 A block (which is interpreted as a point mass) on a frictionless table attached to two

stretched springs of rest length l0 with spring constant k.(a) Motion along the x

direction and (b) motion along the y direction.

Solution

There are two ways to solve this problem, and it worth looking at both methods. In the first

method, we calculate the potential energy U of springs as a function of x and y and then

calculate forces by dierentiation. In the second, we calculate force directly, remembering

that the force is always along the direction of the stretched spring. The energy of the two

spring system is

1

U= k(|l1 l0 |2 + |l2 l0 |2 ),

2

where l1 and l2 are the lengths of springs 1 and 2, respectively.

(a) For motion along the x-direction, l1 = 2l0 + x and l2 = 2l0 x. Then

1

U= k[(2l0 + x l0 )2 + (2l0 x l0 )2 ] = k[l02 + x2 ],

2

p

F x = dU/dx = 2kx = m x, which gives an oscillation frequency of !0 = 2k/m,

p

which corresponds to a period of T = 2/!0 = 2m/k. Alternatively F x can be

calculated directly. The first string pulls in the direction that will tend to restore it to its

rest length l0 , i.e. in the minus-x direction. Its length is 2l0 + x (which is greater than l0

when x > 0), and it exerts a force F1 = k(2l0 + x l0 ). The second string pulls in the

positive-x direction to restore its rest length. The second string has length 2l0 x, and

it exerts a force F2 = k(2l0 x l0 ). The total force from the springs is F1 + F2 = 2kx.

p

(b) For motion along the y-direction, l1 = l2 = (2l0 )2 + y2 .

s

q !2

2 2

y

l1 = l2 = l = (2l0 ) + y = 2l0 1 +

2l0

0 2

1

BBB y C

C y2

= 2l0 B@1 + 2 + CCA 2l0 + .

8l0 4l0

where we have used the Taylor series expansion discussed in 1.6 at the end of the

chapter to approximate the square root expression (see Eq. (1.90)). This is valid in the

t

19 Complex exponential notation

2 !2 !4 3

1 2

66

26 1 y 1 y 7777 1

U(y) = k 2(l l0 ) = kl0 461 + 6 + 7 , kl02 + ky2 , (1.55)

2 2 l0 16 l0 5 2

where in the final step we dropped the small term proportional to (y/l0 )4 . Dierenti-

p

ating gives the force: Fy = dU/dy = ky = my and ! = k/m. Alternatively Fy

can be calculated directly. Let T = k(l l0 ) be the magnitude of the force that each of

the springs exerts on the displaced block. These forces are along the direction of the

displaced springs as shown in Fig. 1.10. The component of the combined force of the

two springs in the y direction is Fy = 2T sin 2T y/(2l0 ), where is the angle the

springs make with the x axis. The spring length l is 2l0 plus terms of order y2 . Thus to

obtain a force to linear order in y, we can replace l by 2l0 to yield T = kl0 and Fy = ky

in agreement with the energy calculation.

If the equilibrium length of the springs were their unstretched length l0 rather than 2l0 ,

then the transverse displacements would be (l l0 ) = y2 /(2l0 ) + , and U(y) ky4 /(2l0 ).

This is an anharmonic spring, whose equation of motion is not so easily solved. This

example shows that it is the stretching of the springs, which puts them under tension, that

is responsible for the harmonic restoring force for transverse displacements. As we shall

see, this is essentially the same physics that controls the vibrations of a violin string.

ing with sine and cosine functions is mathematically cumbersome. However, there exists

an easier way to do algebra with sinusoidal functions, a way that involves the exponential

function with complex arguments. This section introduces you to this algebra and its use

in describing oscillations.

We introduce this approachpby considering the Taylor series expansion of exp(i), where

is a real number and i 1:

(i)2 (i)3 (i)4 (i)5

ei = 1 + i + + + + ... (1.56)

2! 3! 4! 5!

2 3 4 5

= 1 + i i + +i ... . (1.57)

2! 3! 4! 5!

Collecting the individual terms in the sum in Eq. (1.57) into their real and imaginary parts,

which corresponds to collecting together the even and odd powers of in the Taylor series,

t

20 Oscillations

t

Fig. 1.11 The complex plane with complex number z represented in (a) cartesian coordinate

representation and (b) polar representation. Arrow shows the circle traced out when

= !t, in which case x = r cos !t and y = r sin !t.

gives

! !

i 2 4 3 5

e = 1 + ... + i + ... . (1.58)

2! 4! 3! 5!

We recognize the series inside the two sets of parentheses as the Taylor series for cosine

and sine, respectively. Thus, Eq. (1.58) yields

Equation (1.59) is one of the most amazing and useful formulas in mathematics.1 It is also

the primary reason for introducing complex numbers into physics.

An alternative derivation of Eq. (1.59) is instructive. The complex exponential obeys the

same rules of dierentiation that real exponentials do. Therefore,

d2 i

e = i2 ei = ei , (1.60)

d2

and ei obeys Eq. (1.5), the same dierential equation as the displacement in a simple har-

monic oscillator, but with !0 = 1. Thus ei must be of the form of Eq. (1.6). To determine

!0

the unknown coefficients a and b, we note that ei |=0 = 1 and dei /d = iei ! i. Thus

a = 1 and b = i and Eq. (1.59) is reproduced.

Before delving into the uses of Eq. (1.59), lets consider the graphical representation of

complex numbers. Any complex number z can be written in terms of its real and imaginary

parts: z = x + iy, where x and y are real numbers. Therefore any complex number can

be plotted in an (x, y) coordinate system known as the complex plane, as shown in Fig.

1.11(a). In this representation of complex numbers, the x coordinate gives the real part of

the complex number z and the y coordinate gives the imaginary part. The vector proceeding

from the origin to the point (x, y) in the complex plane is sometimes

p called a phasor.

We can also represent z in terms of its absolute value r |z| = x2 + y2 and the angle

1 Setting = , Eq. (1.59) becomes, with a little manipulation, ei + 1 = 0, an economical and almost magical

relationship between the two most fundamental transcendental numbers, e and , the multiplicative and

additive identities, 1 and 0, and the fundamental imaginary number, i.

t

21 Complex exponential notation

that it makes with the real (x) axis. Noting that x = r cos and y = r sin , we obtain

z = x + iy (1.61)

= r cos + i r sin (1.62)

i

= re . (1.63)

Equation (1.63) is sometimes called the polar representation of z. Figure 1.11(b) shows

the polar representation of z in graphical form. We can convert between the cartesian and

polar representations of a complex number guided by the diagrams in Fig. 1.11. To convert

from the polar to the cartesian representation z = rei , we use

x = r cos (1.64)

y = r sin . (1.65)

To convert from the Cartesian to the polar representation, we use

p p q

r = zz = (x + iy)(x iy) = x2 + y2 (1.66)

! y

Im z

= tan 1 = tan 1 . (1.67)

Re z x

Figure 1.11 also provides us with a picture of how an oscillating sinusoidal function is

represented by the real part of a complex exponential. If we let the angle = !t, then as

t increases, z = rei!t sweeps out a circle of radius |z| = r moving counterclockwise in the

complex plane, as shown in Fig. 1.11(b), with a sinusoidal projection onto the real axis

given by x = Re(z) = r cos !t.

Armed with this knowledge about complex exponentials, how do we use them to describe

oscillations? There are several dierent options. Below, we introduce two that are widely

used and that we will find useful in our study of oscillations.

We start by pointing out that e i!t is a solution to Eq. (1.5), the harmonic oscillator

equation, provided ! = !0 , as you can readily verify. Thus, there are two independent

solutions, as required, and the general solution can be expressed by

i!0 t

(t) = a1 e + a2 ei!0 t . (1.68)

The coefficients a1 and a2 can in general be complex, which provide four independent

parameters. We can eliminate two of those parameters, by insisting that (t) be real, as we

would expect for physically meaningful solutions. Demanding that (t) be real means that

the imaginary part of (t) is zero. Thus, (t) and its complex conjugate are equal:

i!0 t

a1 e + a2 ei!0 t = a1 ei!0 t + a2 e i!0 t

. (1.69)

Collecting terms together gives

i!0 t

(a1 a2 )e + (a2 a1 )ei!0 t = 0 . (1.70)

i!0 t

Because e and ei!0 t are independent functions, the only way Eq. (1.70) can be zero is if

t

22 Oscillations

the coefficients (a1 a2 ) and (a2 a1 ) are zero. Therefore, requiring that (t) be real means

that a2 = a1 (or equivalently a1 = a2 ). This leaves only two independent parameters, the

amplitude and phase of a1 , a satisfying situation since we know that there are only two free

parameters in the general solutions. Expressing a1 = (A/2)ei in terms of its real amplitude

A/2 and its phase leads to

(t) = A cos(!0 t ), (1.71)

where we have used a2 = a1 = (A/2)e i . Thus we see that Eq. (1.68) provides a completely

general solution equivalent to Eq. (1.10). The amplitude A and phase , or equivalently the

real and imaginary parts of a1 (= a2 ) are determined by the initial conditions.

Starting from the form given by Eq. (1.10) (or Eq. (1.71)), we can also write the general

solution to Eq. (1.5) as

h i

(t) Re Ae i!0 t , (1.72)

where we have introduced the complex amplitude A = Aei . In this case, we use only one

of the two complex solutions, the real part, and thus have only two constants, the amplitude

A and phase of A.

We will find the complex form of Eq. (1.72) to be quite useful for the treatment of

cases when there is damping and when there are periodic external driving forces. There are

many other equivalent ways of expressing (t) in terms of a complex function. For example,

replacing !0 and by !0 and or replacing by /2 and taking the imaginary part

of the resulting expression provide representations that are equivalent to Eq. (1.72). We

will almost exclusively use Eq. (1.72). We will also find the complex form of Eq. (1.68) to

be useful in certain situations.

Exercise 1.5.1 The motion of a certain oscillator is described by the real part of the

equation

B

x(t) = ei!t (1.73)

15 + 2i

where B is a real number. What are the amplitude and phase of x(t)? That is, rewriting

Eq. (1.73) in the form x(t) = A cos(!t ), find expressions for the amplitude A and

phase , where A and are real numbers. Hint: find the magnitude r and the phase

as illustrated in Fig. 1.11.

As another example, suppose you want to express cos(+ ) in terms of sines and cosines

of the angles and and not their sum. You could rack your memory, look this up, or you

could simply derive the correct expression using complex exponentials. Following the last

route, we first express cos( + ) in terms of complex exponentials, remembering that we

are keeping only the real parts

cos( + ) = ei(+ ) = ei ei (1.74)

= (cos + i sin )(cos + i sin ) (1.75)

= (cos cos sin sin ) + (1.76)

i (sin cos + cos sin ) . (1.77)

t

23 Series expansions

Keeping only the real parts at the end of our calculation, we arrive at the purportedly

familiar formula

cos( + ) = cos cos sin sin . (1.78)

Using the complex exponential notation, there is no need to remember or look up such

formulas as they are readily derived in a few lines. In doing the calculation, however, it

is important to keep all the terms until the end to make sure that when you discard the

imaginary parts, you are not inadvertently discarding some real part of the expression.

Exercise 1.5.2 Following the same procedure shown above, find a similar formula for

sin( + ) in terms of sines and cosines of the angles and . Hint: Explain why you

should start by writing sin( + ) = i ei(+ ) .

Exercise 1.5.3 Starting from the expression ei = cos + i sin , show that

1 i i

sin = 2i (e e ) (1.79)

1 i i

cos = 2 (e + e ). (1.80)

Exercise 1.5.4 The hyperbolic sine and hyperbolic cosine functions are defined as

sinh x = 1

2 ex e x

(1.81)

1 x x

cosh x = 2 e +e . (1.82)

Exercise 1.5.5 Consider the function

h i

y(t) = Re a1 ei!t + a2 e i!t , (1.83)

where a1 and a2 are complex constants. Show that an alternative way of writing y(t)

is

where A1 and A2 are real constants. In particular, derive expressions for A1 and A2

in terms of a1 and a2 . Using Hint: Write a1 = 1 i 1 and a2 = 2 + i 2 where 1 ,

1 , 2 , and 2 are all real. Expand the complex exponentials in Eq. (1.83) using Eq.

(1.59) and then find expressions for A1 and A2 in terms of and .

We have made extensive use of series expansions in this chapter, usually in order to sim-

plify the mathematical analysis of a problem. All of them can be derived from the general

expression for a Taylor series expansion of a function f (x):

f 00 (a) f 000 (a)

f (x) = f (a) + f 0 (a)(x a) + (x a)2 + (x a)3 + ... (1.85)

2! 3!

t

24 Oscillations

Here we list series expansions for a few common functions. They can be derived from Eq.

(1.85) but as they are used so often, you should commit them to memory:

ex = 1 + x + 1 2 1 3 1 4

2! x + 3! x + 4! x + ... (1.86)

1 3 1 5

sin x = x 3! x + 5! x ... (1.87)

1 2 1 4

cos x = 1 2! x + 4! x ... (1.88)

1 2 1 3 1 4

ln(1 + x) = x 2x + 3x 4 x + ... (1.89)

n

(1 + x) = 1 + nx + 2! x + n(n 1)(n

n(n 1) 2

3!

2) 3

x + ... (1.90)

The Taylor expansion for (1 + x)n , known as the binomial expansion, is often overlooked

by new physics students. Nevertheless, it is one of the most useful expansions in physics. It

works for both positive and negative values of n, as well as for both integer and non-integer

values.

p

Example 1.3 Find the Taylor series expansion of f (y) = 1/ a2 y2 .

Solution

We start by rewriting f (y) in the form (1 + y)n

1 1

f (y) = p = p (1.91)

a2 y2 a 1 (y/a)2

" y 2 # 1/2

1

= 1 . (1.92)

a a

Equation (1.92) is now in the form of Eq. (1.90) where x = (y/a)2 and n = 1/2. Thus,

the Taylor series expansion is

2 ! 2 ! ! ! y 2 !2 3

1 6666 1 y 1 1 3 77

f (y) = 641 + + + ...7775 (1.93)

a 2 a 2 2 2 a

" #

1 1 y 2 3 y 4

= 1+ + + ... (1.94)

a 2 a 8 a

Frequently, only the the first two terms are needed when using the binomial expansion.

Thus, it is handy to memorize it: (1 + x)n ' 1 + nx for |x| 1.

Exercise 1.6.1 Find the first two non-vanishing terms in the Taylor series for the fol-

lowing expressions (the first terms are given; you should fill in the ..." term). To do

so, first cast the relevant part of each expression in the form (1 + x)n , where |x| 1,

and then use Eq. (1.90) to perform the expansion. The expansion should be a poly-

nomial in the expansion variable similar to the example above.

t

25 Summary of important points of Chapter 1

1 Q

E x (x) = p (1.95)

40 x x2 + a2

1 Q

' (1 + ...) where x a

40 x2

1

K(3) = ( 1)mc2 , where = p (1.96)

1 32 /c2

1 2

' m3 (1 + ...) v c

2

Equation (1.95) gives the electric field E x along the x axis generated by a uniformly

charged nonconducting wire of length 2a oriented along the y axis and carrying a

total charge Q. Equation (1.96) gives the relativistic expression for the kinetic energy.

Does each expression reduce to its expected result in the limits specified above?

Briefly explain (one sentence each!).

At the end of this and every chapter, we provide a summary of the chapters main points.

You should understand all of the points summarized below. The descriptions in the sum-

mary are kept brief deliberately. Refer back to the chapter for more extensive explanations

and for illustrative examples.

Oscillations can occur when there is a minimum in the potential energy and the system is

displaced from that minimum. The restoring force, given by the gradient of the potential

energy, pushes the system towards equilibrium in the vicinity of the minimum.

The restoring force is usually not linear, or equivalently, the potential energy is not a

parabola. A nonlinear force law leads to a nonlinear equation of motion, which in general

does not have an analytic solution (exceptions do exist, however). The first line of attack

for such a problem is to expand the force law in a Taylor series about the equilibrium

position and to keep only the linear term. The resulting equation of motion is linear and

has analytical solutions. The linearized force law can also be obtained by expanding the

potential energy in a Taylor series about the equilibrium position and keeping terms that

are quadratic in the displacement from the equilibrium position.

Whenever you linearize a force law, you must take care to determine the range of dis-

placements over which the linearized force law provides an accurate description of the

system. An estimate for the displacement at which the error becomes bigger than the

desired accuracy (e.g. 1% or 10%) can be determined by keeping the next non-vanishing

term in the Taylor series. It is important to understand that some force laws can have a

relatively large range of displacements over which the linear force law is valid, while

other force law may have very small linear ranges. The linearized force law for pendu-

lum, for example, is accurate over a relatively large range of displacements, about one

t

26 Oscillations

third of a half revolution, while the linearized Morse potential for a diatomic molecule

is valid over displacements only on the order of 1% of the distance between atoms.

The oscillatory motion that results from a linearized force law is called simple harmonic

motion.

Obtaining solutions to problems involving oscillations about a linearized force law are

greatly simplified by using complex exponentials instead of sine and cosine functions.

Because the linearized equation of motion describing free oscillations is a linear homo-

geneous second-order dierential equation, it has two independent solutions, sine and

cosine, and the most general solution is a linear combination of these two solutions. The

general solution also has two integration constants, which, for a specific problem, can be

determined from the initial conditions. Specifying the initial conditions usually means

specifying the initial position and velocity of the system. If external forces, like a grav-

itational field, are added, the equation of motions becomes inhomogeneous. For simple

oscillators such as those treated here, the inhomogeneous equation can be converted into

a homogeneous equation by choosing the coordinates so that the equilibrium position of

the system corresponds to zero displacement.

t

27 Problems

Problems

1.1 A particle undergoing simple harmonic motion travels a total distance of 6.98 cm

in a full period T = 1.71 s. (a) What is the average speed of the particle? (ans:

4.08 cm/sec) (b) What are its maximum speed and acceleration? (ans: 6.41 cm/s,

23.6 cm/s2 )

1.2 A mass m = 0.2 kg attached to the end of a spring with spring constant k = 1 N/m is

released from rest at t = 0 s from an extended position xm . After 0.5 s, the speed of

the mass is measured to be 1.5 m/s. Calculate xm , the maximum speed, and the total

energy. (Ans: xm = 0.75 m, vm = 1.67 m/s, E = 0.28 J)

1.3 A volume V of air is contained in a vertical cylinder of radius R that is closed at the

bottom and fitted with piston of mass m at the top. The piston fits snugly so that no

air escapes the cylinder. Nevertheless, assume that the piston moves freely up and

down without friction. The equilibrium height of the bottom of the piston is h such

that the equilibrium volume of air is R2 h.

(a) When the piston is displaced a distance x from its equilibrium position, show

that the restoring force is given by:

" ! #

h

F(x) = mg 1 (1.97)

h+x

where = C p /Cv = 1.4 for air. Assume that the piston is displaced from its

equilibrium position for a short enough time that no heat enters or leaves the

cylinder of gas.

(b) Expand the force law in a Taylor series about the pistons equilibrium position

to quadratic order in the displacement x. What is the eective spring constant of

this system?

(c) Find an expression for the oscillation frequency of the piston for small displace-

ments about its equilibrium position. If the mass of the piston is 10 kg and the

equilibrium height is 0.10 m, what is the period of oscillation? Based on your

answer, do you think our assumption that there is no heat flow in or out of the

cylinder is justified?

1.4 In the text, we showed that the reduced mass = m/2 for a diatomic molecule made

of two identical atoms each of mass m. Show that the reduced mass of a diatomic

molecule made up of atoms of unequal masses m1 and m2 is given by = m1 m2 /(m1 +

m2 ). Verify that this reduces to the above result for the case that m1 = m2 . To do this,

follow the same methodology used in the text. Write down equations of motion for

each atom separately and then subtract one from the other to obtain an equation of

motion for the dierence R2 R1 , or for R2 R1 r0 . Hint: After obtaining the

equations for R1 and R2 separately, divide through by the appropriate mass before

subtracting the two equations to obtain an equation of motion for d2 (R2 R1 )/dt2 .

t

28 Oscillations

Lennard-Jones2 potential

" 12 6 #

U(r) = 4 , (1.98)

r r

where and have units of length and energy, respectively, with values that depend

on the particular atoms involved.

(a) Plot the Lennard-Jones potential as a function of r. More specifically, plot U(r)/

as a function of r/ for the range 0.9 r/ 3. Scale the vertical axis so the

energy well is clearly visible.

(b) Show that the minimum energy for the Lennard-Jones potential is located at r =

r0 21/6 . Physically, what does this distance represent? What is the minimum

energy? What is the binding energy for this potential?

(c) Show that the spring constant for the Lennard-Jones

p potential is k ' 57.1/ 2

and the oscillation frequency by ! ' 7.56 // where is the reduced mass

for the system.

(d) On the same plot that you plotted the Lennard-Jones potential, plot its harmonic

approximation, which includes all the terms in the Taylor expansion of U(r)

about r = r0 up to and including the term quadratic in r r0 .

(e) The vibrational frequency f of a nitrogen molecule is about 8.2 1013 Hz. Using

numbers found on the web or in books, determine the values of and . Using

those values, use the results from part (c) to find the the width of the well ?

What is the value of the ratio /r0 ?

1.6 (*) A particle of mass m moves in a potential

(a) What is its equilibrium position x0 (nearest to x = 0)? Express you answer in

terms of V0 , and mg. You may assume that mg < V0 . (ans: x0 = (1/) sin 1 (mg/V0 ))

(b) What is the frequency of small oscillations about the equilibrium position. Your

answer may be expressed as a function of A, , m, and x0 . (ans: ! = [(V0 /m) cos x0 ]1/2 =

(/m)1/2 [(2 V02 m2 g2 ]1/4 )

1.7 A block of mass m = 0.1 kg attached to an ideal spring moves on a frictionless

surface. Let x(t) be its displacement as a function of time.

(a) At time t = 0, the block passes through the point x = 0 moving to the right (i.e.,

toward positive x). At time t = 0.5 s, the block reaches its maximum excursion

of xm = 10 cm.

(i) What is the period T of oscillation? (Ans: t = 2 s)

(ii) What is the spring constant k? (Ans: k = 0.12 N/m)

(iii) What is the maximum velocity 3m ? (Ans: 0.314 m/s)

(iv) Write the full expression for x(t).

2 Sir John Lennard-Jones (18941954) was a professor of theoretical chemistry at Cambridge University. He

was one of the pioneers of molecular orbital theory.

t

29 Problems

t

Fig. 1.12 (a) Block at rest prior to time t = 0. (b) Configuration of the block at t = 0+ (i.e., right

after t = 0) shown in lighter gray and configuration for general time t in darker gray.

p

(b) At time t1 = 5T/8, a particle of mass m and moving with velocity 3 = 33m / 2 to

the right collides completely inelastically with the block (i.e., collides and sticks

to the block) of the oscillator in (a).

(i) What is the position of the block as a function of time after the

p collision?

p p

2 1

(Ans: x(t) = (3/2)xm cos[ T (t t1 ) + ] where = tan 2).

(ii) At what time does the block pass the origin again? (Ans: 1.53 s).

1.8 A block of mass m, resting in its equilibrium position on a frictionless platform, is

attached to a Hookes-law spring with spring constant k with one end fixed on a

vertical wall connected to the platform as shown in Fig. 1.12. At time t = 0, the

platform is suddenly tilted upward to an angle .

(a) Determine the position x(t) of the block for all t > 0. (ans: x(t) = l0 (mg/k) sin (1

p

cos !0 t), where !0 = k/m)

(b) What are the total kinetic energy K(t) and potential energy U(t) for times t > 0.

At what times is the kinetic energy a maximum? (ans: t = (1 + 2n)(T/4), where

T is the period of oscillation)

1.9 A cylindrical cork of radius R = 1.5 cm, height d = 0.75 cm, and specific gravity of

0.24 floats at an air-water interface.

(a) How far below the surface of the water is the bottom of the cork in equilibrium:

(ans: h = 0.18 cm).

(b) At time t = 0, a small of aluminum bead of mass mA = 1.0 gm is placed on

top of the cork while imparting a downward velocity of 30 = 10 cm/s to the

combined cork and bead system. Neglecting any friction or water flow induced

by the motion of the cork, calculate the depth y(t) of the bottom of the cork for

all t > 0 and sketch its form. What is the maximum depth and at what time is

it first reached? (ans: 0.55 cm at time 0.04s) What is the minimum depth and at

what time is it first reached? (ans: 0.09 cm at time 0.098s).

1.10 (*) Four identical springs of rest length l0 are attached to a mass m at right angles

to each other on a frictionless horizontal plane as shown below. The two springs

parallel to the x-axis are stretched to length 2l0 while the two springs parallel to the

y-axis have their rest length l0 . When the mass is attached to any one of the individual

springs alone, it oscillates with angular frequency !0 = 2/T 0 . Express your answers

in terms of !0

t

30 Oscillations

t

Fig. 1.13 Figure for problem 1.10

(a) Calculate

p the frequency ! x of small oscillations in the x direction. (Ans: ! =

2!0 )

(b) p

Calculate the frequency !y of small oscillations in the y direction. (Ans: ! =

3!0 )

(c) At time t = 0, the block is given velocities 3 x = 30 and 3y = 30 in the x and y

directions, respectively. What are x(t) and y(t) for t > 0.

1.11 A mass m is attached to two walls a distance l0 away by massless linear springs

whose unstretched length is l0 as shown below in gray. The force of gravity pulls the

mass down and stretches the springs as shown below in black.

t

Fig. 1.14 Figure for problem 1.11

(a) Determine the equation for the equilibrium position y0 of the mass. (Ans: y0 is

the solution to the equation, T l(y) sin = k(l(y) l0 )y, where l(y) = (l02 + y2 )1/2 ).)

(b) What is the frequency of oscillation of the mass about its equilibrium position?

Your answer can be expressed in terms of l(y0 ). (Ans: !2 = k(y0 )/m, where

k(y0 ) = k[l3 (y0 ) l03 ]/l3 (y0 ).)

(c) The same arrangement of springs and masses is placed on a horizontal friction-

less plane. Gravity does not operate. What is the potential energy of displacement

to lowest order in y? What is the equation of motion of the mass to lowest order

in y? (Ans: my = 2ky3 /l02 )

2 Damped Oscillations

Real mechanical oscillations, like those of the pendulum or mass on a spring, die out with

time. This is because there is always some source of friction or damping in macroscopic

systems that robs the oscillator of energy. The damping of oscillations can be undesirable,

as for the case of a pendulum in a clock, or they can be purposely built into the design of a

system, as for a swinging door or the suspension of an automobile. In the case of a clocks

pendulum, the aim is relatively simple: minimize damping. In other cases, however, the

damping needs to be set just right. The swinging door needs to swing, but its oscillations

should die out quickly. If the damping is too weak, the swinging of the door can surprise,

and harm, the next person coming through. If the damping is too strong, the door may

take too long to close or be too hard to open. Studying the physics of damped oscillations

helps us design oscillators that have the right degree of damping required for any particular

application.

We begin our discussion of damped oscillations with the ideal spring we studied in Chap-

ter 1, but this time we include the eect of friction, which slows the pendulum and causes

its oscillations to die away. Friction can arise from air resistance or from friction associ-

ated with the pivot point. Experiments reveal that these sources of friction can usually be

accounted for quite well by a damping force Fd that is proportional to the instantaneous

velocity 3 of the mass:

Fd = 3, (2.1)

In the case of low-Reynolds number viscous friction, = 6R. The minus sign indicates

that the damping force acts in the direction opposite to the velocity of the mass ( > 0), so

its eect is always to slow the motion of the pendulum.

We obtain the equation of motion for the damped harmonic oscillator following the same

procedure we used in Chapter 1, but now including the frictional force given by Eq. (2.1),

d2 x dx

m = kx , (2.2)

dt2 dt

or, dividing through by m, we obtain the linearized equation of motion for the damped

31

t

32 Damped Oscillations

t

Fig. 2.1 Simple pendulum, captured at an instant when it is moving downwards towards = 0.

The damping force Fd = v is in the opposite direction of the velocity v.

harmonic oscillator:

d2 x dx

+ + !20 = 0 , (2.3)

dt2 dt

p

where !0 k/m, and = /m. For the case of no damping where = 0, solutions to

Eq. (2.3) are sinusoidal. More precisely, the solutions are given by x = A cos(!0 t ), so

we see that !0 is just the oscillation frequency of the undamped oscillator. The amplitude

A and the phase are constants determined by the initial conditions.

When the damping constant is small, we expect the solutions still to be approximately

sinusoidal but with an amplitude that gets smaller with the passage of time. Based on these

physical considerations, a plausible guess for the solution would be a function of the form

t

x(t) A e cos(!t ). (2.4)

Equation (2.4) is an oscillating solution with an amplitude a exp( t) that decays with

time. The constant must have units of inverse time so that the argument t of the expo-

nential is dimensionless. We might guess that is proportional to because has units of

inverse time and, more importantly, because increasing increases the damping force (see

Eq. (2.1)), which should cause the oscillations to decay more rapidly.

Our guess, based on physical grounds, of Eq. (2.4) as a solution for the damped pen-

dulum turns out to be correct, as you can readily check by substituting Eq. (2.4) into Eq.

(2.3). Doing so yields explicit expressions for and !: = /2 and

q q

! !1 = !20 ( /2)2 = !0 1 ( /2!0 )2 , (2.5)

which is correct in the limit of weak damping, i.e. when 1 ( /2!0 )2 > 0. However, the

solution fails in the limit of strong damping, i.e. when 1 ( /2!0 )2 < 0. This approach can

be salvaged, but the method is mathematically cumbersome, primarily because it involves

sine and cosine functions.

We need a more systematic procedure for finding solutions to Eq. (2.3). While dierent

t

33 The damped pendulum

approaches are available, we shall choose one that is useful not just for the problem at

hand, but one that finds much wider application in other problems, specifically when we

consider more complex oscillators and waves. As we expect the solutions to be oscillatory,

we look for solutions of the form

i!t

x(t) = a e (2.6)

where a and ! can be complex numbers. By allowing ! to be complex, Eq. (2.6) can

express both sinusoidally oscillating and exponentially damped solutions, depending on

the value of !. For example, suppose we let ! = b ic, where t, b, and c are real positive

numbers. Then,

i!t i(b ic)t ct ibt ct

e =e =e e =e (cos bt i sin bt) , (2.7)

ibt

e = cos bt i sin bt . (2.8)

Equation (2.7) is a decaying sinusoidal function, as promised. Its true that the function

is complex, but this turns out not to be a problem. As discussed in 1.5, we shall use the

complex functions of the form ei in order to simplify the analysis, but in the end will keep

only the real part of the solution. The imaginary part comes along for the ride, but is not

used.

The amplitude a in Eq. (2.6) can also be complex, which allows us to adjust not only

the real amplitude of x(t) but also the phase of its oscillations. To see this, we write the

complex amplitude a in polar form as discussed in 1.5

a = Aei , (2.9)

where A and are real numbers. Letting ! = b ic as before, the expression for x(t) in Eq.

(2.6) becomes

i!t

x(t) = a e = Aei e ct

e ibt

= Ae ct

e i(bt )

ct

= Ae [cos(bt ) i sin(bt )] (2.10)

Taking only the real part of Eq. (2.10), (t) = Ae ct cos(bt ), we see that writing Eq.

(2.6) with a complex amplitude a and complex frequency ! allows us complete freedom to

express an exponentially decaying sinusoidal function with arbitrary (real) amplitude and

phase.

Armed with Eq. (2.6), we return to our task of solving the equation of motion for the

damped harmonic oscillator. To this end, we propose Eq. (2.6) as a general solution to Eq.

(2.3). Our task is to determine which sets of values of the constants a and ! in Eq. (2.6)

yield meaningful solutions to Eq. (2.3). Writing the time derivatives as x dx/dt and

x d2 x/dt2 , Eq. (2.3), the linearized equation of motion, becomes

x + x + !20 x = 0 , (2.11)

The form of our proposed solution is x(t) = a exp( i!t), which gives x = i!a exp( i!t)

t

34 Damped Oscillations

and x = !2 a exp( i!t). Substituting these expressions into Eq. (2.11) and then canceling

the common factors of a exp( i!t) yields the equation

!2 i ! + !20 = 0 . (2.12)

The values of ! satisfying Eq. (2.12) are obtained from the quadratic formula

r 2

! = i !20 . (2.13)

2 2

This result leads to three dierent solutions depending on whether the value of the dis-

criminant !20 ( /2)2 is positive, negative, or zero. These correspond to dierent degrees

of damping: underdamped, overdamped, and critically damped. We treat each of the three

cases in the sections that follow.

We begin our analysis with the case when ( /2)2 < !20 , that is when the rate of damping

is smaller than the rate of oscillations. In this limit, you might expect the oscillations to

proceed more or less unscathed, but to slowly die out over time. As we shall see, this is

indeed the case, at least in the extreme underdamped limit when ( /2)2 !20 .

When ( /2)2 < !20 , we can rewrite Eq. (2.13) as

! = i /2 !u , (2.14)

where !u is a positive real frequency given by

q

!u = !20 ( /2)2 > 0 . (2.15)

Equation (2.14) actually implies that there are two independent solutions, one for the

plus sign and another for the minus sign in Eq. (2.14):

i!+ t t/2 i!u t

xu1 (t) = a1 e = a1 e e (2.16)

i! t t/2 i!u t

xu2 (t) = a2 e = a2 e e (2.17)

Aside from the amplitudes a1 and a2 , the two solutions dier only by a minus sign in the

complex exponential, but that is enough to make the two solutions independent, meaning

that xu1 (t) , Cxu2 (t) for any (complex) constant constant C. Note also that xu1 (t) , xu2 ( t).

Our next step is to exploit one of the most important properties of the solutions to linear

dierential equations, namely, the principle of superposition. The principle of superposi-

tion states that if there are two (or more) solutions to a linear dierential equation, then

any linear combination, i.e. any sum of two or more solutions, is also a solution to the

dierential equation in question. Thus, because Eqs. (2.16) and (2.17) are both solutions

to the linearized equation of motion, Eq. (2.3), so is their sum. Therefore, the most general

solution to Eq. (2.3) for this problem is

1 1

xu (t) = a1 e t/2 ei!u t + a2 e t/2 e i!u t

(2.18)

2 2

1 t/2

= e a1 e + a2 e i!u t .

i!u t

(2.19)

2

t

35 The damped pendulum

This is the general solution with coefficients a1 and a2 with four independent parameters

that are still undetermined. As in the case of complex solutions for the undamped oscillator

discussed in 1.5.2, requiring that our solution xu (t) be real means that a1 = a2 = Ae i ,

where A and are real. Substituting into Eq. (2.19) gives

1

xu (t) = Ae t/2 ei(!u t ) + e i(!u t ) (2.20)

2

= Ae t/2 cos(!u t ) (2.21)

This equation can, of course, be expressed in terms of sine and cosine functions:

t/2

xu (t) = e (A1 cos !u t + A2 sin !u t) , (2.22)

where A1 = A cos and A2 = A sin . Equation (2.21) or alternatively Eq. (2.22) express

the general solution to the damped pendulum oscillator

As in the case of the undamped oscillator, the initial conditions determine the two un-

known coefficients (A and or A1 and A2 ). As a specific example, lets consider the case

where the oscillator is initially at rest, but has an initial velocity 30 . We apply these initial

conditions to the general solution given by Eq. (2.22)

xu (0) = A1 = 0 (2.23)

xu (0) = 30 = !u A2 , (2.24)

which gives

30

xu (t) = e t/2 sin !u t.

!u

3i e t/2 q

= q sin !20 ( /2)2 t . (2.25)

2 2

!0 ( /2)

The solution given by Eq. (2.25) is a damped sine wave with a frequency !u given by Eq.

(2.15). It is plotted in Fig. 2.2. Note than !u is smaller than the natural frequency !0 of the

undamped pendulum. This result is consistent with what you might have expected at the

outset, namely that weak damping causes the oscillations to die o with time and also to

slow down somewhat.

There are four (at least!) equivalent ways that the general solution for the underdamped

oscillator can be written. We summarize them here:

8 t/2

>

>

> e (a1 ei!u t + a1 e i!u t )

>

>

>

> t/2

<e

> (A1 cos !u t + A2 sin !u t)

xu (t) = >

> t/2

. (2.26)

>

>

> B e cos(! u t )

>

>

>

:C e t/2 sin(! t )

u

amplitude that decays exponentially in time. Each function has two constants that are de-

termined by the initial conditions for any given problem. In each case, the eect of the

constants is simply to adjust the phase of the sinusoidal oscillations within the decaying

exponential envelope. To solve a problem with any given set of initial conditions, you are

free to choose any of the above functional forms.

t

36 Damped Oscillations

q (t)

1.0 Q = 10

Q=3

Q=1

0.5

0.0 t/T

2 4 6 8

-0.5

t

-1.0

Fig. 2.2 Transient response of an underdamped oscillator for different values of Q. The gray

line shows the decaying exponential envelope of the Q = 10 trajectory.

Oscillators are often characterized by a quantity known as the quality factor, defined here

as

!0

Q= . (2.27)

For the moment, we are considering the case when ( /2)2 < !20 , which means that Q > 12 .

The quality factor Q also roughly corresponds to the characteristic number of oscillations

a free oscillator undergoes before the oscillations eectively die out. Therefore, if a pen-

dulum oscillates 20 times back and forth before the amplitude of its oscillations diminish

appreciably, then Q for that pendulum is about 20 (to within a factor of 2 or so). Figure 2.2

shows the trajectories of a damped oscillator for several dierent values of Q.

From Eq. (2.15) it is also easy to show that the oscillation frequency for an underdamped

pendulum is given by !u = !0 [1 (2Q) 2 ]1/2 . In the limit that Q 1, this means that the

change in the oscillation frequency from the undamped value is ! = !u !0 ' (8Q2 ) 1 ,

which is very small indeed when Q 0. Thus, a weakly damped pendulum oscillates at

very nearly the natural frequency.

When ( /2)2 < !20 , the system is said to be underdamped. This means simply that the

system will oscillate back and forth at least once, and usually many more times. For the

opposite case, when ( /2)2 > !20 , the system does not oscillate but simply relaxes directly

to its equilibrium position. In that case, the system is said to be overdamped. We examine

that case in the next section.

Exercise 2.1.1 For the underdamped case, show that the amplitude of the oscillations is

damped by about a factor of 20 after Q oscillations (i.e. for t = QT , where T = 2/!0

is the period).

t

37 The damped pendulum

q (t)

0.4

Q = 0.5

Q = 0.3

0.2

Q = 0.1

t

0.0 t/T

1 2 3 4

Fig. 2.3 Transient response of an overdamped oscillator for different values of Q. Note the

difference in time scale and the smaller-amplitude response compared to the

underdamped cases shown in Fig. 2.2. The critically damped case, which

corresponds to Q = 1/2, is the most rapidly decaying solution.

For the overdamped case, the discriminant !20 ( /2)2 in Eq. (2.13) is negative, which

means that the frequency ! is purely imaginary. To make this explicit, we rewrite Eq.

(2.13) as

0 r 1

BBB 2 CCC

B

! = i B@ !20 CCA (2.28)

2 2

= i ( /2 ) , (2.29)

where

q

= ( /2)2 !20 > 0 . (2.30)

As in the underdamped case, Eq. (2.29) implies that there are two solutions. Once again,

the most general solution is a linear combination of the two solutions:

( 12 + )t ( 12 )t

xo (t) = b1 e + b2 e (2.31)

t/2 t t

=e b1 e + b2 e , (2.32)

where b1 and b2 are constants that are determined by the initial conditions. Note that in this

case, there is no need for complex coefficients b1 and b2 .

Once again, as an example, lets work out the solution for the initial conditions where

the pendulum starts at = 0 with an angular velocity x(0) = 30 . According to Eq. (2.31),

t

38 Damped Oscillations

xo (t = 0) = b1 + b2 = 0 (2.33)

xo (t = 0) = ( 12 + )b1 ( 12 )b2 = 30 . (2.34)

Solving Eqs. (2.33) and (2.34) simultaneously gives b1 = b2 = i /2 . Substituting this

into Eq. (2.31) gives1

30

xo (t) = e t/2 e t e t (2.35)

2

30

= e t/2 sinh t (2.36)

q !

3i

= q e t/2 sinh ( /2)2 !20 t (2.37)

2 2

( /2) !0

Note that this solution never crosses zero; it merely rises from zero to some maximum and

then decays back down to zero. Such a system is said to be overdamped. The condition for

the overdamped case is ( /2)2 > !20 , which corresponds to Q < 12 . Figure 2.3 shows the

trajectories of overdamped as well as underdamped oscillators. Note that while increasing

the amount of damping (by increasing ) leads to a smaller maximum displacement for the

dierent cases in Fig. 2.3, it also has the eect of slowing the approach to equilibrium. The

critically damped pendulum reaches equilibrium the fastest.

Exercise 2.1.2 You can show that the general solutions of the overdamped and under-

damped cases are equivalent. Begin by noting that !u = i . Then show that setting

!u in Eq. (2.25) equal to i yields Eq. (2.36). Thus, the solutions to the underdamped

and overdamped cases are obtained from the a single set of equations.

Exercise 2.1.3 For t 2/ , show that in the strongly overdamped limit, meaning that

Q 1, Eq. (2.37) can be written as

30 2

x(t) ' e Q t . (2.38)

the response of the of a damped oscillator in the strongly overdamped limit. Hint:

Keep only the positive exponential in the sinh function (why is this justified?) and

then expand the argument of the exponential using the binomial expansion (see Eq.

(1.90)).

For the critically damped case, the discriminant !20 ( /2)2 = 0 and ! in Eq. (2.13) is

given by

!= i (2.39)

2

1 Here we use the hyperbolic sine function, denoted sinh, and defined as sinh x 12 (e x e x ). There are

corresponding hyperbolic cosine and tangent functions defined as cosh x 12 (e x + e x ) and

tanh x = sinh x/ cosh x.

t

39 The damped pendulum

From the form of our proposed solution given in Eq. (2.6), one may be temped to write the

solutions as

t/2

xc (t) = c e , (2.40)

and be done with it. But there is a problem with Eq. (2.40): It cannot describe the solution to

the situation where the initial displacement xc (t = 0) is zero and the initial angular velocity

xc (t = 0) is finite (i.e. not zero). Another clue that there is a problem comes from the fact

that only one integration constant, C, appears in Eq. (2.40). As a solution to a second order

dierential equation, we expect two integration constants.

So how do we obtain the correct solution for the critically damped case? There are

several approaches available but perhaps the most obvious is to use either the underdamped

or overdamped solutions we have already found, and to take the limit of those solutions

as the discriminant ( /2)2 !20 goes to zero. After all, if we make only an infinitesimal

change in the relative values of and !0 , we expect only an infinitesimal change in the

solution, so such a procedure should lead to a well-defined result. In fact, this procedure

works very well. Problem 2.2 guides you through the mathematics. The result is

t/2

xc (t) = (c1 + c2 t) e , (2.41)

where c1 and c2 are the two (expected) integration constants whose values are fixed by the

initial conditions.

Once again, as an example, lets work out the solution for the initial conditions where

the pendulum is initially at x(0) = 0 with an angular velocity of x(0) = 30 . In this case, the

solution is quite straightforward: (0) = 0 implies that c1 = 0 in Eq. (2.41), while x(0) = 30

implies that c2 = 30 . Thus, the specific solution for these initial conditions is

t/2

xc (t) = 30 t e . (2.42)

Starting at t = 0 with an initial velocity of i , the pendulum moves away from its equilib-

rium position, according to Eq. (2.42), and reaches a peak displacement time t = 2/ , after

which the displacement decays exponentially towards zero.

In Fig. 2.4, we summarize the three cases of damping treated above: underdamped, over-

damped, and critically damped. Here we show the case where each oscillator is started

at rest from a finite initial displacement. While derived in the context of the the damped

pendulum, these results are completely general and apply to all linear simple harmonic

oscillators. Notice that only the underdamped case exhibits oscillations. Figure 2.4 also

shows that an overdamped oscillator initially decays more rapidly than a critically damped

oscillator, but ultimately decays more slowly. This is an important consideration when de-

signing dampers for door as well as many other practical devices where oscillations are

undesirable. For these initial conditions, adjusting the damping constant so that it cor-

responds precisely to the case of critical damping ensures that the system returns to the

equilibrium position the fastest.

t

40 Damped Oscillations

q (t)

1.0

critically damped

0.5

overdamped

underdamped

0.0 t/T

1 2 3 4

-0.5

t

Fig. 2.4 Three cases of damped oscillations: underdamped (solid), overdamped (short

dashes), and critically damped (long dashes).

The amplitude of the oscillations of a damped pendulum decreases with each successive

swing as energy is lost to the friction of the damping force. The energy lost is simply the

work done by the damping force on the pendulum

Z

W= Fd ds , (2.43)

where the minus sign arises according to Eq. (2.1) Fd = 3 = m3. Substituting this

into Eq. (2.43) and using the fact that ds = 3 dt gives the work W(t) done by the damping

force since time t = 0

Z t

W(t) = m 32 (t0 ) dt0 . (2.44)

0

The work-energy theorem tells us that the change in energy is the work done on the system

by outside forces, in this case the friction force. Thus

Z t

E(t) E0 = W(t) = m 32 (t0 ) dt0 . (2.45)

0

dE

= m 32 (t) . (2.46)

dt

Thus, the instantaneous rate of energy loss is proportional to 32 .

Of course the total energy can also be calculated directly from the instantaneous sum of

t

41 Energy loss from damping

1 2 1

E(t) = K(t) + U(t) = m3 (t) + kx2 (t) . (2.47)

2 2

Calculating E(t) is fairly tedious for the most general case, but simplifies considerably in

the limit of very weak or very strong damping, that is, in the strongly underdamped and

overdamped regimes. It also becomes a bit easier to calculate the results if we consider

specific initial conditions. So we use the initial conditions introduced previously of zero

initial displacement and a finite initial velocity. In the very weak damping limit, which

corresponds to ( /2)2 !20 or Q 1, we can use Eq. (2.25) and set !u ' !0 . This gives

30 t/2

x(t) = e sin !0 t. (2.48)

!0

Dierentiating this expression, we obtain the velocity

30 h i

3(t) = l x(t) ' e t/2 !0 cos !0 t 1

2 sin !0 t . (2.49)

!0

Because /2 !0 , we can neglect the sine term above. Thus

t/2

3(t) = x(t) ' 30 e cos !0 t . (2.50)

1 2 1

E(t) = m3 (t) + kx2 (t) (2.51)

2 2

1 2 t 1

' m30 e cos2 !0 t + m320 e t

sin2 !0 t (2.52)

2 2

1 2 t

= m30 e (2.53)

2

1

= m320 e !0 t/Q . (2.54)

2

Thus we see that in the very weak damping limit, energy decays exponentially in time with

a time constant of 1/ . Or equivalently, since Q = !0 / , energy decays on a scale that is

Q/2 times slower than the period of an oscillation.

In the strong damping limit, which corresponds to ( /2)2 !20 or Q 1, we start with

Eq. (2.36)

30 t/2

x(t) = = e sinh t (2.55)

30

= e t/2 e t e t (2.56)

2

30 h ( /2)t i

= e e ( + /2)t . (2.57)

2

Except for an initial short transient, exp( ( + /2)t) exp(( /2)t). Therefore, to

determine the long-time behavior, we need to keep only the first exponential in Eq. (2.57).

We can further simplify matters by expanding the quantity ( /2), keeping in mind that

t

42 Damped Oscillations

( /2)2 !20 :

q

= ( /2)2 !20 (2.58)

2 2

2s 3

666 4!02 777

6

= 6664 1 1 777 (2.59)

2 2 75

20 1 3

666BBB 1 4!20 CC 77

= 4@1 6 B + ...CCA 1775 (2.60)

2 2 2

!20

' (2.61)

where in going from Eq. (2.59) to Eq. (2.60) we have used the binomial expansion (1 x)a =

1 ax + O(x2 ). With these approximations, Eq. (2.57) becomes

30 (!20 / )t

x(t) ' e . (2.62)

2

Dierentiating, we obtain

30 !20 (!20 / )t

3(t) ' e . (2.63)

2

Substituting these two expressions into Eq. (2.47) yields

1 2 1

E(t) = m3 (t) + kx2 (t) (2.64)

2 2

1 2 2 1 2 2(!20 / )t

' m30 Q ( 4 Q + 1) e (2.65)

2

1

' m320 Q2 e 2Q!0 t . (2.66)

2

Thus, in the strong damping limit where Q 1, the energy decays at a rate of 2Q!0 while

in the weak damping limit where Q 1, it decays at a rate of !0 /Q. This suggests that the

energy decays most rapidly in the vicinity of critical damping or Q 12 .

Exercise 2.2.1 A 0.5 kg mass hangs on the end of a 1.0-m-long pendulum that is part

of a grandfather clock. In ten minutes, the amplitude of the pendulum decays from

4.0 cm to 0.5 mm. (a) What is the Q of this pendulum? (b) Show that the fraction of

energy lost each cycle is approximately 2/Q.

Exercise 2.2.2 A church bell produces a sound at 220 Hz (the sound of the musical

note A below middle C). After it is struck, the sound dies down in about 5 s. What is

the approximate Q of the bell?

Exercise 2.2.3 For a very weakly damped oscillating pendulum with amplitude 0 ,

show (i) that the total energy is given by ET ' 12 ml2 !20 02 , (ii) that the energy dis-

sipated in one cycle is W f ' 12 l2 !20 02 (2/!0 ), and (iii) that the ratio of energy

stored to energy lost in one cycle is ET /W f ' !0 /2 = Q/2. This last equation

Q = 2ET /W f is often used as the definition of Q. Its worth remembering.

t

43 Shock absorbers

Perhaps the most common use of damped oscillators is to mitigate the potentially delete-

rious eects of a sudden shock to a system. A sensitive piece of electronics like a mobile

phone is more likely to break if it falls one the floor unprotected than if has an elastic rub-

ber coating to cushion the impact. We dont want the cushion to be too springy, however,

as it is generally preferable for the phone to bounce only once or twice before coming to

rest. The springiness of the protective rubber coating should be appropriately damped.

Another example of how damped oscillators are used is found in the shock absorbers of an

automobile. The idea is for the shock absorbers to cushion the impact to the car (and its

passengers) when it runs over a bump or hole in the road by mounting springs between the

wheels and the car frame. However, we would not want the springs to cause running over

a bump to send the car into large persistent oscillations, which could cause loss of control

(of either the car or the passengers stomachs). Any oscillations that are caused by the sus-

pension springs should die out quickly after the springs have done their job of cushioning

the impact.

Lets consider a simple model of a car suspension system. The suspension system con-

sists of a heavy duty spring and a viscous damper, which we model as a dashpot," as

shown in Fig. 2.5. The spring connects the car frame (green) to axle (black dot) of the

wheel (large gray circle). The dashpot, which acts as a damping mechanism, consists of

(1) hollow cylinder with a closed bottom (red) that is connected to the axle and (2) a piston

(blue) that fits inside the cylinder and is connected directly to the car frame. When the

car frame moves relative to the axle, the spring either stretches or compresses while the

dashpot provides the viscous friction force proportional to the velocity of piston (and the

car frame) relative to the cylinder (and the wheel axle). The distance from the axle to the

car frame is z. We take z = z0 to be the resting distance between the car frame and the axle.

t

Fig. 2.5 Suspension system with shock absorber on an automobile wheel. Here, the dashpot

(red and blue) fits inside the hollow of the spring (black).

t

44 Damped Oscillations

In this state the net force F on the car is zero: F = k(z0 zu ) Mg = 0, where k is the

spring constant, zu corresponds to the position of the unstretched spring, and M is the mass

of the car that is supported by this spring. When the car moves vertically relative to the

axle, the viscous damping force provided by the dashpot is given by Fd = z, where z is

the vertical velocity of the car and is the viscous damping coefficient, which is taken to

be constant.

Putting this all together, the equation of motion for the vertical movement of the car is

F = k(z zu ) Mg z = Mz . (2.67)

Defining = z z0 as the distance from the equilibrium position z0 and using the equilib-

rium condition that k(zu z0 ) = Mg, Eq. (2.67) becomes

k = M . (2.68)

Rearranging terms, this becomes

+ + !20 = 0 (2.69)

p

where = /M and !20 = k/M. This is the same equation of motion we previously

encountered for the damped pendulum (see Eq. (2.3)). Notice the trick we used above to

get rid of the constant (inhomogeneous) term involving zu in Eq. (2.67). By defining the

zero of the coordinate to be at the equilibrium position of the spring under the weight of

the car, and using the equilibrium condition that k(zu z0 ) = Mg, we obtained Eq. (2.69),

an equation of motion without a constant term. This is the same trick we introduced in

1.3.2. Its a good trick to remember!

Lets use our results to design the suspension for a car. We suppose that one wheel

supports about 400 kg or about 1/4 of the mass of a 1600 kg car. Lets suppose that the car

is moving at about 3 x = 15 m/s (54 kph ' 34 mph) when it runs over a sudden drop in the

road of d = 5 cm. For the sake of simplicity we shall assume that the distance between the

car frame and the axle remains fixed at z0 during the fall. We further assume that when the

tire hits the ground it abruptly stops its vertical motion while the car frame initially keeps

moving. Its initial vertical velocity is 3z = (2gd)1/2 ' 1.0 m/s downward. Thus, the initial

conditions are = 1.0 m/s and = 0 (i.e. z = z0 ).

To get a sense of what is going on, we start by assuming that we want the suspension

to be critically damped. The boundary conditions at t = 0 are zero displacement and a

nonzero velocity xii . Our equations of motion are identical to those for the damped spring,

Eq. (2.42), with (t) replacing x(t):

c (t) = i t e t/2

. (2.70)

In designing the shock absorber, there are at least two important considerations. First, we

do not want the shock of the car hitting the ground after the drop to be too large. Second,

we want the car to recover quickly from the eects of the drop in the road. Lets consider

the shock first. Our concern is simply that the vertical acceleration associated with the car

hitting the ground not be too big. Dierentiating Eq. (2.70) twice to obtain the vertical

acceleration of the car reveals that the maximum acceleration amax occurs at t = 0 and

is given by amax = i , that is, the damping constant times the initial vertical velocity

t

45 Shock absorbers

when the wheel his the ground. Since this is a modest size drop, lets see if we can limit

amax to 12 g ' 5 m/s2 . (For reference, the accelerations experienced while walking have an

amplitude of about 15 g.) With = 1.0 m/s, this gives = 5 s 1 .

We also want the car response and recovery time to occur on a time scale no longer

than a second or so. There are two time scales in the problem, the damping time, which

is given by 2/ , and the inverse frequency (or period) given by !0 1 = 1/Q . We expect

the optimal solution to be near critical damping, for which Q = 12 . Thus, both time scales

2/ and !0 1 = 1/Q are on the order of 0.4 s. Since the maximum acceleration is set by

the damping constant, we will hold fixed and vary !0 to vary Q near its value for the

critically damped case.

In Fig. 2.6 we plot the response given by Eq. (2.70). For the critically damped case,

the spring compresses ( < 0), reaches a maximum compression of about 15 cm at about

t = 0.4 s, and then relaxes back to zero displacement after about 3 s.

To get a better sense of how to optimize our design, we examine what happens for the

slightly underdamped and overdamped cases. Once again we can adapt our results from

the damped oscillator to get the responses for the underdamped and overdamped cases.

From Eqs. (2.25) and (2.36), we obtain the following results for the underdamped and

overdamped cases respectively

i t/2

u (t) = e sin !u t (2.71)

!u

i t/2

o (t) = e sinh t (2.72)

where !u is given by Eq. (2.15) and is given by Eq. (2.30). In Fig. 2.6. Returning to our

design criterion that the maximum acceleration be given by amax 5 m/s2 , we once again

calculate amax from Eqs. (2.25) and (2.36) and find that amax = i for both the under-

damped and overdamped cases, just as we found for the critically damped case. Therefore,

we shall keep the same value of = 5 s 1 . For the underdamped and overdamped cases,

there is another time scale that we can tune, namely !0 1 . For the critically damped case, !0

is fixed since Q = !0 / = 12 , but for the other cases it can be varied since Q = !0 / > 12

for the underdamped case and Q = !0 / < 12 for the overdamped case. Because we expect

that optimal damping occurs near critical damping, we vary Q by only a factor of two in

either direction choosing Q = 1 and Q = 14 for the underdamped and overdamped cases,

respectively.

The acceleration profiles for the three cases are very similar, with the critically damped

and overdamped cases being nearly indistinguishable (see bottom of Fig. 2.6). On the other

hand, the displacement curves for the three cases dier significantly from each other. For

the overdamped case, the wheel is still more than 5 cm from its equilibrium position 4 s

after impact. While such a long relaxation time might have been expected, it is interesting

that the wheel returns to its equilibrium position considerably more rapidly for the under-

damped than for the critically damped case. By increasing Q while keeping constant for

the underdamped case, we have increased the spring constant k, since k = M!20 = M 2 Q2 .

The increased value of k limits the maximum displacement, which explains why the maxi-

mum displacement is smallest for the underdamped case. Of course, the underdamping also

t

46 Damped Oscillations

causes the overshoot as the wheel approaches it equilibrium position. Increasing Q > 1

would increase the number of oscillations and thus defeat the damping function of the

shock absorber. Nevertheless, the curves in Fig. 2.6 suggest that a small degree of un-

derdamping, which implies a stier spring, is desirable in the design of shock absorbers.

However, it is clear that the price paid for returning to equilibrium quicker is a somewhat

stier ride, as indicated by somewhat greater average absolute acceleration evident in Fig.

2.6. This trade o between a softer ride but slower response and a stier ride but quicker

response is one reason why the family sedan has a more comfortable ride than a sports car,

but the sports car gives the driver more control.

For our final design, we therefore set = 5 s 1 and Q = 1. For M = 400 kg, this

gives a spring constant of k = M!20 = M 2 Q2 = 10, 000 N/m and a damping constant of

= M = 2, 000 N s/m. Real cars have significantly stier springs. One reason is that

0.0 t (s)

1 2 3 4

Q=1

Q = 1/2

0.1

Q = 1/4

0.2

a(t) (m/s2 )

5

0 t (s)

1 2 3 4

t

Fig. 2.6 Response of automobile shock absorbers to running over a pothole. Top:

displacement from equilibrium wheel position. Bottom: vertical acceleration of

automobile. The damping is the same for each curve; Q is varied by changing the

spring constant k = M 2 Q2 .

t

47 Shock absorbers

the air-filled tires take up much of the shock of driving over a bump, so the shock absorbers

only have to do only a part of the job.

t

48 Damped Oscillations

Problems

2.1 Consider a pendulum consisting of a mass m = 0.20 kg at the end of a massless rod

1.20 m long attached to the ceiling. The mass is displaced 5.0 cm from equilibrium,

let go, and allowed to oscillate freely. After 10 oscillations, its amplitude is 4.0 cm.

(a) Is the system overdamped, underdamped, or critically damped? Is the value of Q

much greater than 1 or much less than 1 (do not calculate Q yet). Briefly explain

(1 or 2 short sentences!).

(b) Using the fact that the amplitude of oscillations decays approximately as e t/2 ,

determine the approximate numerical value (correct to within 5%) of Q for this

system.

2.2 In this problem you will show that the general solution for the critically damped

oscillator is given by Eq. (2.41), which for your convenience we rewrite here

t/2

c (t) = (c1 + c2 t) e . (2.73)

Start with the solution for the overdamped case, Eq. (2.31), written below in a slightly

more convenient form

h i

(t) = e t/2 b1 e t + b2 e t , (2.74)

q

where = ( /2)2 !20 . The strategy is to find (t) from Eq. (2.74) in the limit

! 0, which corresponds to the critically damped case. Since you can make

arbitrarily small, you can expand the exponentials e t and e t in a power series about

t = 0. Go ahead and do this. Then gather together all the even and odd terms of the

power two power series and show that can you obtain

" !

t/2 ( t)2 ( t)4

(t) = e (b1 + b2 ) 1 + + + ...

2! 4!

!#

( t)2 ( t)4

+ (b2 b1 ) t 1 + + + ... . (2.75)

3! 5!

Letting c1 = (b1 + b2 ) and c2 = (b2 b1 ) , you have the solution given by Eq. (2.41).

How is it that the terms in the power series go to zero as ! 0 but the term (b2 b1 )

can survive and remain finite?

2.3 Find the solutions for the damped pendulum for the initial conditions that x(t = 0) =

= 0) = 0. In particular,

x0 > 0 and (t

(a) starting from Eq. (2.19), show that

!

t/2

xu (t) = x0 e cos !u t + sin !u t (2.76)

2!u

for the underdamped case.

t

49 Problems

!

t/2

xo (t) = x0 e cosh t + sinh t (2.77)

2

for the overdamped case.

(c) Starting from Eq. (2.41), show that

t t/2

xc (t) = x0 1 + e (2.78)

2

for the critically damped case.

(d) On a single graph, plot of x(t) vs. t for Q = { 14 , 12 , 1, 2, 4} and x0 = 0.05 over

the range 0 t 5. For the purpose of your plots set the period of oscillation

T = 1 ) !0 = 2.

(e) Which of your plots are overdamped, underdamped, or critically damped. For

any underdamped cases, how much is the period changed from what it would

have been if there were no damping?

2.4 According to classical electrodynamics, an accelerating charge radiates energy at a

rate of Kq2 a2 /c3 where a is its acceleration, q is the charge, c is the speed of light,

and K = 6.0 109 N m2 /C2 . Suppose that the electron is oscillating during during

one cycle of its motion according to x = A sin !t.

(a) Show that the energy radiated during one cycle is given by Ke2 !3 A2 /c3 , where

e is the electronic charge.

(b) Recalling that the energy of a simple harmonic oscillator is 12 m!2 A2 , show that

Q = mc3 /Ke2 !.

(c) Using a typical value of ! for visible light, estimate the lifetime" of the radiating

system.

2.5 Redo the problem in 2.3 of a car going over a drop in the road but relax the assump-

tion that the spring does not extend while the car is falling. Use the same parameters

used in 2.3. Take the radius and mass of the car wheel to be 0.225 m and 8 kg, re-

spectively. Determine the time-dependent displacement and acceleration. Present

your results in graphical format similar to that of Fig. 2.6.

2.6 A circular metal plate with moment of inertia I = 0.02 kg-m2 is suspended from a

metal wire. When the plate is twisted through an angle away from its equilibrium

position, it experiences a torque k. If the plate rotates between two magnet heads,

it experiences a dissipative torque b. Thus the equation of motion for the plate in

the presence of both the magnets and an external torque ext (t) is

I + k + b = ext (t) .

When the magnet heads are removed, there is no dissipative torque, and the system

undergoes simple harmonic motion with period T = 0.5 s. When the magnet heads

are in place, the system has a quality factor Q = 8.

(a) What is the torque constant k? Be sure to specify units. (Ans: k = 0.322 kg-m2 /s2 ))

(b) What is the damping constant ? Be sure to specifyunits. (Ans: = 0.01 kg-m2 /s)

t

50 Damped Oscillations

t

Fig. 2.7 Torsion Oscillator

(c) At time t = 0, the plate is given an impulse such that its initial angular velocity

is radians/sec. What is the angular displacement after 0.125 seconds?

2.7 A mass-and-spring harmonic oscillator oscillates at a frequency !0 = 2/T 0 when it

is undamped. It is placed in a viscous medium, and it acquires a quality factor Q = 1.

(a) At time t = 0 the mass, at rest at its equilibrium position at the origin, is de-

livered at impulse that gives it a velocity 30 . Determine the time T 1 at which

the mass first returns to the origin. What is the position x2 of the mass p at time

T 1 /2?. Express

p your answer

p in terms of T 0 and 30 . (Ans: T 1 = T 0 / 3, x2 =

[30 T 0 /( 3)] exp[ /(2 3])

(b) The oscillator is immersed in another viscous medium, and its quality factor

decreases to 1/4. If the particle is displaced a distance x0 from the origin at time

t = 0 and released, what is its position at time T 0 ? Ans:

x0

x(T 0 ) = p + e 2 e 2+ ,

2 3

p

where = 2 3

2.8 (French 3-14) An object of mass 0.2 kg is hung from a spring whose spring constant

k is 80 N/m. The object is subject to a resistive force given by 3, where 3 is the

velocity in m/s.

(a) Set up the dierential equation

p of motion for the free oscillation of the system.

(b) If the damped frequency is 3/2 of the undamped frequency, what is the value

of the constant ? [Ans: = 4 kg/s]

(c) What is the Q of the system, and by what factor is the amplitude ofposcillation

reduced after 10 complete cycles? [Ans: Q = 1, A(10)/A(0) = e 20/ 3 = 1.76

10 16 .]

2.9 (French 3-15) Many oscillatory systems, although the loss or dissipation mechanism

is not analogous to viscous damping, show an exponential decrease in their stored

average energy with time E = E 0 e t . A Q for such oscillators may be defined by

using the definition Q = !0 / , where !0 is the natural angular frequency.

(a) When the note middle C" on the piano is struck, its energy of oscillation de-

creases to one half its initial value in about 1 sec. The frequency of middle C is

256 Hz. What is the Q of the system? [Ans: Q = 2/ ln 2 = 2321]

t

51 Problems

(b) If the note an octave higher (512 Hz) takes about the same time for its energy to

decay, what is its Q. [Ans: Q = 4642]

(c) A free, damped harmonic oscillator, consisting of a mass of m = 0.1 kg moving

in a viscous fluid of damping coefficient (Fviscous = 3), and attached to

a spring of spring constant k = 0.9 N/m, is observed as it performs damped

oscillatory motion. Its average energy decays to 1/e of its initial value in about

4 sec. What is the Q of the oscillator? What is the value of ? [Ans: Q = 12,

= 0.025 kg/s]

2.10 Consider a damped harmonic oscillator with a natural frequency !0 and inverse

damping time whose displacement as a function of time is given by x(t).

(a) Derive an expression for x(t) for the over damped case subject to initial condi-

tions x(0) = x0 and x(0) 30 .

(b) Derive the corresponding expressions for x(t) for the critically damped and under-

damped oscillators using only the expression for the over-damped oscillator

and

q the facts that = 0 in the critically damped oscillator and ! i!1 =

i !20 ( /2)2 in the under-damped oscillator. Verify that x(0) = x0 and 3(0) =

v0 in both cases.

3 Resonance

good pendulum clock typically has an accuracy of about 1 part in 105 , or about 10 seconds

a month. However, a lone pendulum without any other mechanism is a poor time-keeping

device because the oscillations that measure the passage of time gradually die out as fric-

tion saps energy from each successive swing. The only way to keep a such a clock going

is to periodically inject some energy. This is how the pendulum on a mechanical clock

sustains its motion. By means of a mechanism involving a ratchet and a falling weight or

a compressed spring, the pendulum is given a small impulse each cycle. A key feature that

permits the impulse to be so small is that it is synchronous with the natural frequency of the

pendulum. That is, the impulse has the same frequency !0 as the natural frequency of the

pendulum. This phenomenon, that periodic forcing near the natural frequency of an oscil-

lator results in oscillations with much larger amplitudes than forcing at other frequencies,

is called resonance.

Resonance is widely observed in nature and exploited by technology. One of the princi-

pal uses of resonance is the detection of waves, most notably sound and electromagnetic

waves, including radio, light, and x-rays. For example, radio waves, a vast spectrum of

which fills almost every space we occupy every hour of every day, can be detected when

they set into motion electrons in an antenna. The radio waves create currents that oscillate

over an incredibly broad range of frequencies, from several hundred kilohertz to hundreds

of megahertz. How is it that the radio tuner finds the signal of that one radio station you

want to listen to? The answer lies in the phenomenon of resonance.

The aim of this chapter is to provide you with a quantitative introduction to oscillations

and resonance. By the end of the chapter, you should understand how a radio receiver

tunes in a specific frequency broadcast, how a seismometer detects a broad range of seis-

mic movements, and how a pendulum works. And so we proceed, from the simple to the

sublime.

As noted above, the only way to maintain the motion of an oscillator indefinitely is to

periodically inject some energy. One example is a pendulum clock, which was previously

discussed. Another example is a child on a swing with a parent periodically giving the

child a push. The pushing is most eective if it occurs in sync with the natural period of

52

t

53 Forced oscillators and resonance

the swing and if it is in the proper direction. That is, the parent should push the swing in

the direction it is already going, and certainly not in the opposite direction, which would

slow down the swing. For the forcing to be most eective, therefore, it should occur in

sync with the natural frequency of the oscillator and it must be properly phased. We will

look for these features when we obtain solutions to the equation of motion for a forced

oscillator.

3.1.1 Resonance

Here we consider a weight of mass m suspended by a spring of spring constant k and sub-

jected to a friction force proportional the weights velocity 3, Fd = 3, which can arise,

for example, from a surrounding viscous fluid or from confinement to a hollow cylinder

providing lubricated contact between the weight and the cylinder as depicted in Fig. 3.1.

The piston is driven by a sinusoidally oscillating force F(t) = F0 cos !t, which might be

provided by coupling magnetically to a metallic weight.

Measurement of the time-dependent position z(t) of the weight as a function of ! and

F0 reveals the following properties:

1 After an initial period after the driving force is turned on, the system reaches a steady

state in which the vertical position of the weight undergoes simple harmonic motion at

the same frequency ! as the driving force but with a phase lag, (!), relative to the

driving force: z(t) = z0 + A(!) cos(!t ), where z0 = mg/k is the equilibrium position

of the weight and A(!) is the maximum amplitude.

p

2 When ! !0 = k/m, the driving force and the weight displacement follow each other

( (!) 0): they reach maxima and minima at nearly the same time. The amplitude A(!)

increases with increasing !.

3 When ! !0 , the driving force and displacement are almost completely out of phase,

: the driving force is at a maximum when the displacement is at a maximum and

vice versa. The displacement amplitude decreases with increasing frequency.

t

Fig. 3.1 Mass in a hollow cylinder suspended by a spring and driven by an oscillating

magnetic field from coil. The motion is damped by the lubricated contact between the

mass and the cylinder.

t

54 Resonance

4 The amplitude A(!) reaches a maximum in the vicinity of ! = !0 with a peak value that

is higher the smaller the damping (reduced for example by switching from a viscous

fluid like glycerine to a less viscous one like water). The peak also gets narrower as

damping is decreased. This is resonance. The phase displacement-force phase lag is /2

at ! = !0 : displacement is zero when the driving force is a maximum.

5 The power required to drive the system follows a curve as a function of ! that is similar

to that of the amplitude with a peak at !0 and dying o with ! at large !, but going to

zero at ! = 0.

We now show how all of these properties follow from the harmonic-oscillator equation

in the presence of a sinusoidal driving force. Gravity acts downward on the weight in

the geometry of Fig. 3.1. Following the procedure discussed in 1.3.2, we define x(t) =

z(t) mg/k to remove gravity from the equation of motion, to obtain

d2 x dx

m = kx + F0 cos !t . (3.1)

dt2 dt

After some rearrangement of terms, this equation can be rewritten as

d2 x dx F0

2

+ + !20 z = cos !t , (3.2)

dt dt m

p

where !0 k/m and /m. Equation (3.2) has the same form as Eq. (2.3), the

equation for the unforced damped oscillator, except that Eq. (3.2) includes a forcing term,

which is an explicit function of the independent variable t. Thus, Eq. (3.2) is inhomoge-

neous, whereas Eq. (2.3) is homogeneous, as discussed in 1.1.

We look for solutions to Eq. (3.2) using the same approach we used for the unforced

pendulum, by assuming a solution of the form x(t) = Ae i!t , where A is complex. To take

advantage of the algebraic simplicity that exponential notation aords, we write the forcing

term as

i!t

F0 cos !t = F0 Re e , (3.3)

i!t

taking only the real part of e , which is cos !t. Thus, Eq. (3.2) becomes

x + x + !20 x = (F0 /m) e i!t

, (3.4)

where we have omitted explicitly writing Re for notational simplicity. We just have to

remember to take only the real part of any solution we obtain at the end of our calculation.

Initially, we looking only for steady state solutions, that is, for solutions that have only

a simple sinusoidal time dependence of the form

i!0 t

x(t) = Ae , (3.5)

where we take the weight to oscillate at a frequency !0 to allow for the possibility that the

weight oscillates at a frequency dierent from either the natural frequency !0 or the drive

0

frequency !. Keep in mind that A and e i! t are complex, which allows for there to be a

phase dierence between the driving force and the response of the oscillator. In the end we

will keep only the real part of x(t). The other terms in the equation, including !0 , , m, and

F0 are real constants that describe the oscillator and the nature of the forcing.

t

55 Forced oscillators and resonance

To look for solutions, we substitute Eq. (3.5) into Eq. (3.4), which gives

!0 2 i !0 + !20 A e i! t = (F0 /m) e i!t .

0

(3.6)

0

Multiplying both sides of the equation by ei! t yields

!0 2 i !0 + !20 A = (F0 /m) e i(! !0 )t

. (3.7)

Notice that all the time dependence in this equation is on the right hand side. The left

hand side is completely independent of time! For the two sides of the equation to be equal,

the right hand side must be independent of time as well. The only way this can happen is

if !0 = !. Thus, even before finding a full solution to the problem, we arrive at a very

important result: If a linear oscillator is driven at a frequency !, then the oscillator will

respond only at that frequency. This means that it wont respond at the natural frequency

!0 . Now this turns out to be true only in steady state, that is, for times much longer than

the natural damping time of the oscillator, which from the previous section is typically of

the order of 1 = Q/!0 , where Q is the quality factor introduced in 2.1.1. It turns out

that we have left out the transient solution in our work above. Dont worry, we will put it

in shortly. Nevertheless, the point is important enough to reiterate. If a linear oscillator is

driven at a frequency !, its steady state response will be only at that frequency.

With !0 = !, our trial solution becomes

i!t

x(t) = A(!)e , (3.8)

where we display explicitly that A depends on frequency. We will not always do so in what

follows. Proceeding with this trial solution, Eq. (3.7) becomes

!2 i ! + !20 A = F0 /m . (3.9)

F0 /m

A(!) = . (3.10)

!20 !2 i !

Writing the amplitude in polar form

, (3.11)

where A0 (!) and (!) are the real amplitude and phase, Eq. (3.10) yields 1

F0 /m

A0 (!) = q 2 (3.12)

!20 !2 + 2 !2

0 1

BB ! CCC

(!) = tan 1 BB@ CA . (3.13)

!20 !2

i!t i(!t )

x(t) = Ae = A0 e . (3.14)

p

1 Its worth pausing here to prove to yourself that 1/(a ib) = ei / a2 + b2 , where tan = b/a.

t

56 Resonance

A0 (w)

15

Q = 16

10

Dw 3g Q 1

0 w/w0

0 1 2 3

f (w)

p

Q = 16

p/2 Q = 1/2

t

0 w/w0

0 1 2 3

Fig. 3.2 Amplitude (top) and phase of oscillator response z(t) to periodic forcing of f0 cos !t for

values of the quality factor Q from the lightest to the darkest curves: 12 , 1, 2, 4, 8, 16.

The height of the peak of the resonance curve (top) is given by a0 (!0 ) = Q( f0 /!20 ) and

width ! is proportional to 1/Q in the limit that Q 1. Phase (bottom) of the oscillator

response for values of Q ranging from 1 to 16, with shadings corresponding to the top

curves. The curves corresponding to Q = 1 are shown as dashed lines.

F0 /m

= q 2 cos(!t ), (3.15b)

!20 !2 + 2 !2

where the phase is given by Eq. (3.13). Equation (3.15) represents the complete solution

to the steady state driven damped oscillator. This important result states that amplitude and

the phase of the oscillator are functions of the the drive frequency !, the natural frequency

!0 , and the damping rate .

Equation (3.15) provides the steady state driven solution for x(t) when the physical driv-

ing force is F0 cos !t. What happens when the physical driving force is F0 sin !t or even

F0 cos(!t ) for arbitrary angle ? One only has to replace !t by !t , with = /2

for the F0 sin !t case, in x(t). Thus if for = /2, x(t) = A sin(!t ).

t

57 Forced oscillators and resonance

Figure 3.2 shows the amplitude and phase of the oscillator given by Eqs. (3.12) and

(3.13). The most striking feature of the plots is the dramatic increase of the amplitude A0

of the oscillator when the driving frequency ! is near the natural oscillation frequency !0

for the case of weak damping, that is, when Q 1. This dramatic increase in the response

of the oscillator for driving frequencies near !0 is known as resonance and is a universal

feature of high-Q oscillators.

The frequency at which the amplitude is a maximum is given by

s

1

!max = !0 1 , (3.16)

2Q2

which, for Q 1, is very nearly equation to the free (or natural) frequency of oscillation

!0 , which is often referred to as the resonant frequency.

For Q 1, the maximum amplitude of the oscillator is given by

F0 F0

A0 (!max ) ' A0 (!0 ) = Q =Q , (3.17)

m!20 k

which means that for Q 1, the peak amplitude is enhanced by a factor of Q over the the

displacement F0 /k that would be expected for a constant (non-oscillatory) force F0 .

The range of frequencies around the resonant frequency where the amplitude is large

is determined by the damping

p rate . The amplitude is reduced to half its peak value at

frequencies ! ' !0 ( 3/2) . The so-called full width at half maximum (fwhm), which

is defined as the width ! of the resonant curve at half of its maximum amplitude A(!0 ) is

given by

p !0

! 3 , (3.18)

Q

The quality factor Q thus provides a convenient quantitative measure of the resonant curve:

the peak height increases in proportional to Q and the width narrows in proportion to Q 1 .

Exercise 3.1.1 Starting from Eq. (3.12), show that the for Q 1, the full width half

maximum frequency is given by Eq. (3.18). Hint: In the large Q limit where !0 ,

you can approximate the terms in the denominator of Eq. (3.12) using

!20 !2 = (!0 !)(!0 + !) 12 ! (2!0 ) = ! !0 (3.19)

! !0 . (3.20)

Justify these and any other approximations you make in arriving at Eq. (3.18).

The phase (!) of the oscillator, defined in Eq. (3.15a) and given by Eq. (3.13), also

varies with the frequency !, especially in the vicinity of ! = !0 . For frequencies much

smaller than the resonant frequency, is nearly zero, meaning that the response of the

oscillator is in phase with the forcing. This makes sense: if you push very slowly on a os-

cillator, the applied force is simply balanced by the spring restoring force. Indeed, because

the acceleration and velocity are very small in this case, the first two terms in Eq. (3.4) are

very small compared to the third (see Eq. (3.9)). Thus, the equation of motion becomes

x(t) ' (F0 /m!20 ) cos !t = (F0 /k) cos !t, which is consistent with Eq. (3.15b) in the limit

t

58 Resonance

z(t) F(t)

5 w = 12 w0

0 t/T0

1 2 3

5

z(t) F(t)

5

w = w0

0 t/T0

1 2 3

5

z(t) F(t)

5

w = 2w0

0 t/T0

1 2 3

5

t

Fig. 3.3 Motion z(t) of oscillator (black lines) in response to periodic forcing F(t) (dashed gray

lines) for (a) ! = !0 /2 (below resonance), (b) ! = !0 (at resonance), and (c) ! = 2!0

(above resonance). In all cases, the forcing is given by F(t) = F0 cos !t. Time is

measured in units of the natural period T 0 = 2/!0 and Q = 8 for all curves.

that ! !0 . In this case the amplitude is just the equilibrium displacement F0 /k asso-

ciated with a static force of F0 . The driving force and oscillator response in this limit are

illustrated in Fig. 3.3(a).

As the frequency ! approaches the resonant frequency !0 , (!) increases and the re-

sponse of the oscillator starts to lag behind the driving force. That is, the oscillator reaches

its peak a short time after the forcing goes through its peak value. Precisely at resonance

when ! = !0 , (!0 ) = /2 and the oscillator position lags the driving phase of the driv-

ing force by 90 . On the other hand, the velocity x is in phase with the driving force at

resonance. Thus, the force is always pushing in a direction that tends to enhance the oscil-

lators motion. This is consistent with your experience of pushing someone on a swing; you

push them in the same direction they are already moving. The driving force and oscillator

response at resonance are illustrated in Fig. 3.3(b).

As the driving frequency increases, the oscillator can no longer keep up with the driving

force. When this occurs, the driving force actually opposes the motion of the oscillator for

part of its cycle, resulting in a loss in amplitude. For ! !0 , the driving force is 90 out

t

59 Forced oscillators and resonance

c

15 Q = 16

10

5 c0

c 00

0 w /w0

1 2 3

15

t

10

0 00

Fig. 3.4 Response functions of an oscillator: in phase (solid line) and out-of-phase

(dashed line).

of phase with the velocity so that the force actually opposes and enhances the oscillator

motion for nearly equal amounts of time resulting in a vanishingly small amplitude. In this

case, the displacement and driving force are almost completely out of phase. This situation

is illustrated in Fig. 3.3(c).

We end this section by summarizing the response of the driven oscillator, both in ampli-

tude and phase, at the resonant frequency !0 and for frequencies much less than and much

greater than !0 :

8 8

>

>

> (F0 /m!20 ) = F0 /k for ! ! 0 >

>

>0 for ! ! 0

>

>

< >

>

<

A0 (!) ' >

> Q (F0 /m!20 ) = QF0 /k for ! ! !0 , (!) ' >

>/2 for ! ! !0 . (3.21)

>

>

> >

>

>

:(F0 /m!2 ) ! 0 for ! ! 1 : for ! ! 1

an external periodic driving force in terms of the amplitude and phase of the oscillator.

An alternate way to express this same information about the response of the system is in

terms of the in-phase and out-of-phase amplitudes of the motion. One reason why this

is interesting is that there is a close relationship between the energy dissipated by the

oscillator and its out-of-phase amplitude.

For an oscillating driving force of F(t) = F0 cos !t, we found we found the solution

x(t) = Ae i!t , where A and e i!t are both complex. Writing the amplitude in terms of its

t

60 Resonance

real and imaginary parts A = AR + iAI , the displacement of the oscillator becomes

n o

x(t) = Re A e i!t (3.22)

= Re {(AR + iAI ) (cos !t i sin !t)} (3.23)

= AR cos !t + AI sin !t , (3.24)

where we have kept only the real part of the expression. The cos !t term is, of course, in

phase and the sin !t is out of phase with the driving force F0 cos !t. Thus, from Eq. (3.24)

we see that amplitudes of the in-phase and out-of-phase response are given by AR and

AI the real and imaginary parts of A. With this in mind we define a frequency-dependent

response function for each part

0 AR (!) 1 !20 !2

(!) = = (3.25)

F0 m (!20 !2 )2 + ( !)2

00 AI (!) 1 !

(!) = = 2

. (3.26)

F0 m (!0 ! ) + ( !)2

2 2

A(!) 0 00

(!) = = (!) + i (!) , (3.27)

F0

0 00

where (!) and (!) are its real and imaginary parts. It is clear that

, (3.28)

p

where | (!)| = A0 (!)/F0 = ( 00 )2 + ( 00 )2 and (!) = tan 1 ( 00 / 0 ) is equal to phase of

A of Eq. (3.13).

Figure 3.4 shows the frequency-dependent in-phase and out-of-phase response functions

0

(!) and 00 (!), respectively. The limiting forms of these functions are are summarized

as

8 8

>

>

> 1/k for ! ! 0 >

>

> 0 for ! ! 0

>

>

< >

>

<

0 00

(!) ' >

> 0 for ! ! !0 , (!) ' >

> Q/k for ! ! !0 . (3.29)

>

>

> >

>

>

: 1/(m!2 ) ! 0 for ! ! 1 :0 for ! ! 1

A number of properties of (!) should be noted. 0 (!) is even under ! ! !, positive for

! < !0 , equal to zero at !0 = 0, and negative for ! > !0 . 00 (!), on the other hand, is odd

under ! ! ! and thus equal to zero at ! = 0, positive for all ! > 0, peaked, like A0 (!)

and | (!)|, in the vicinity of ! = !0 with a width that goes to zero with , and equal to Q/k

at ! = !0 . The limiting forms of Eq. (3.29) can be deduced directly from the equations

of motion, Eqs. (3.2) and (3.9). At ! = 0, the inertial (m x) and dissipative ( x) terms are

zero, and all that remains on the right-hand side of the equation is the static term kx, and

(0) = 1/k; when ! ! 1, the inertial term dominates so that (! ! 1) = 1/(m!2 );

finally at ! = !0 , only the dissipative term remains, and (! = !0 ) = i/(m!0 ), implying

that the response is purely imaginary in this limit and that, as a result, = /2.

In the current context, the equations of motion are linear in both x(t) and F(t), and (!)

is simply A(!)/F0 . In most problems, however, A or its analog depend nonlinearly on F(t),

t

61 Forced oscillators and resonance

but the response function (!) is still well defined: it is just the derivative of A with respect

to F0 evaluated at F0 = 0. Thus, (!) is just a generalization to non-zero frequency of a

static susceptibility, like the magnetic susceptibility, @M/@H, where M is the magnetization

and H is the magnetic field, or the electric susceptibility @P/@E, where P is the electric

polarization and E is the electric field. Thus, we expect (! = 0) to be the zero-frequency

susceptibility dx/dF0 = 1/k as it is. As we shall see in a later chapter, (!) is the Fourier

transform of a time-dependent response function (t, t0 ) relating changes in x at time t to

forces at other times t0 : Z 1

x(t) = dt0 (t, t0 )F(t0 ). (3.30)

1

Linear response functions are one of the most important tools for studying properties and

excitations of materials from metals to polymers. Interestingly, the response functions for

these much more complex functions not only share the same underlying mathematics as

our simple example but also exhibit the same behavior with narrow peaks in frequency

near elementary excitations.

3.1.3 Dissipation

It is the periodic forcing that sustains an oscillators motion and keeps it from succumbing

to the damping force. We need to calculate how much work is done by the external force in

one cycle. This energy is in the end lost to friction, which heats the system. Physically, you

can think of this as the energy dissipated by the oscillator. The rate at which the external

force does work is

dW

= F(t)3(t). (3.31)

dt

Pd (w)

Pd Q

15

Q = 16

10

1

2 Pd Dw g Q 1

t

0 w

0 w0

Fig. 3.5 Resonant power absorption as a function of frequency for three different values of Q.

The maximum Pmax and half maximum 12 Pmax as well as full width half maximum

frequency width ! are indicated for the Q = 16 curve.

t

62 Resonance

Recall that in the steady state, both F(t) and 3(t) oscillate at the same frequency ! and

are thus periodic functions of time with period T = 2/!. This allows us to calculate the

average power dissipated in a cycle:

Z Z

1 T 1 T dF

Pd = dtF(t) x(t) = dt x(t), (3.32)

T 0 T 0 dt

where to obtain the second equation, we integrated by parts using the fact that F(T )x(T ) =

F(0)x(0). Using Eq. (3.15a) for x(t) and setting A0 (!) = F0 | (!)|, we obtain

Z

F2

Pd = 0 dt !| (!)| cos !t sin[!t (!)]

T

1 1

= !F02 | (!| sin (!) = F02 ! 00 (!). (3.33)

2 2

It is a straightforward exercise to show that this result is identical the rate that the system

does work against dissipative forces. [See exercise 3.1.3]

Thus, we arrive at the very important result that the rate of energy dissipation is pro-

portional is proportional to ! 00 (!). For this reason, 00 (!) is sometimes referred to as the

dissipation. Note that 00 (!) is odd in ! so that ! 00 (!) is even. In thermodynamic equi-

librium, the power dissipation must be positive (energy is dissipated not created by friction

forces), implying that ! 00 (!) is positive. The positivity of ! 00 (!) in the present case is

associated with the positive sign of . Its sign was chosen so that the viscous or friction

force opposes motion of the oscillator mass. The sign is consistent with the transfer of en-

ergy to incoherent degrees of freedom of the fluid or medium responsible for friction and

to positive power absorption.

Figure 3.5 shows a plot of Pd (!) for three dierent values of Q. In contrast to the

frequency-dependent amplitude A0 (!), the power absorption function Pd (!) is peaked ex-

actly at !0 , and the full width at half maximum is precisely ! = = !0 /Q [Prob. 3.5] for

all , even when > 2!0 and the solutions to the homogenous equation of motion has no

oscillatory part, though the peak can be quite asymmetric about !0 . Setting ! = !0 reveals

that the height of the power resonance curve is Pd = (F02 /2m ) / Q.

Exercise 3.1.2 Using the fact that the force exerted by the oscillating weight on the

medium producing friction is Fd and, thus, the work it does against friction is 3dx,

show that average power loss due to friction is equal to the average power, Eq. (3.33),

power spent by the the external force.

We noted at the outset of this chapter that a well-designed mechanical clock only needs a

very small kick each cycle to keep it going. What this means, of course, is that the energy

stored in the swinging pendulum of a mechanical clock is much greater than the energy

dissipated each cycle. Having already looked at the energy dissipated in a driven oscillator,

lets calculated the energy stored.

The energy of a harmonic oscillator is the sum of its kinetic and potential energies, which

t

63 Forced oscillators and resonance

1 2 1

E= m3 + m!20 x2 , (3.34)

2 2

where the spring constant is given by k = m!20 . Using Eq. (3.15), z = A0 cos(!t ) and

z = !A0 sin(!t ), so that Eq. (3.34) becomes

1 h i

E = mA20 !2 sin2 (!t ) + !20 cos2 (!t ) . (3.35)

2

We are less interested in the instantaneous energy than we are in the average energy. There-

fore we average Eq. (3.35) over one cycle. The only time varying parts are the sine and

cosine terms. Over a single period, sin2 (!t ) and cos2 (!t ) oscillate between 0 and

1

1 with an average value of 2 . Performing these averages, Eq. (3.35) leads to the following

expression for the average steady state energy stored in a driven oscillator

1

E s = ml2 A20 !2 + !20 . (3.36)

4

Using the expression for A0 given by Eq. (3.12), this becomes

F02 !2 + !20

E s (!) = . (3.37)

4m !2 !2 2 +

2 !2

0

To make direct comparison with Pd (!), we integrate Eq. (3.33) over one period to obtain

the energy dissipated in one cycle

Z

2

Ed (!) = Pd dt = T Pd = Pd (!) . (3.38)

cycle !

The ratio of the energy stored to the energy dissipated is given by

Es 1 !2 + !20

= , (3.39)

Ed 4 !

which diverges for ! !20 / and for ! !0 and has a broad minimum at ! = !0 ,

meaning that a little more energy is dissipated relative to the energy stored near resonance

than at other nearby frequencies. At ! = !0 , the ratio E s /Ed reduces to

E s (!0 ) Q

= . (3.40)

Ed (!0 ) 2

This is what we intuitively expect: higher Q oscillators dissipate a smaller fraction of their

energy than do lower Q oscillators. Equation (3.40) is often used as an alternative to Eq.

(2.27) as a definition of Q. Its worth remembering that Q is proportional the energy stored

divided by the energy dissipated in a driven oscillator at resonance.

3.1.5 Transients

periodic force is suddenly applied, starting at time t = 0. Because of its finite inertia, the

oscillator will not immediately start moving at its steady state velocity. Instead, it will take

t

64 Resonance

some time for it to build up its motion. As the oscillator begins to move, it will exhibit some

transient behavior that over time will give way to the steady state oscillations we studied

in the previous sections. From our study of the damping of free oscillations of oscillators

in 2.1, we would expect any transients to die out after approximately Q oscillations, but

as Q can easily be on the order of 104 , it could be quite awhile before a oscillator reaches

its steady state, that is, until its behavior is independent of the initial conditions. Here we

take a look at the transient motion of an oscillator as it builds up to its steady state.

To determine the transient behavior of the oscillator, we once again make use of the

principle of superposition, but in a dierent guise this time. Consider two solutions, x1 (t)

and x2 (t), to the motion of an oscillator: x1 (t) is a solution to Eq. (2.3), the equation of

motion for an unforced free oscillator, and x2 (t) is a solution to Eq. (3.4), the equation of

motion for a forced oscillator:

x1 + x1 + !20 x1 = 0 (3.41)

x2 + x2 + !20 x2 = (F0 /m) cos !t . (3.42)

Suppose we sum these two solutions to create a new function x3 (t) = x1 (t) + x2 (t). Now,

lets substitute x3 (t) into the left hand side of Eq. (3.42) and see what we get:

x3 + x3 + !20 x3 = x1 + x1 + !20 x1 + x2 + x2 + !20 x2 (3.43)

According to Eq. (3.41), the term inside of the first set of parentheses on the right hand

side of Eq. (3.43) is zero, while according to Eq. (3.42), the term inside the second set

of parentheses is F0 cos !t. Therefore x3 (t) is also a solution to Eq. (3.4), the equation of

motion for a forced oscillator. This illustrates an important mathematical result: the sum of

a solution to the homogeneous dierential equation and a of a solution of the inhomoge-

neous dierential equation is also a solution of the inhomogeneous dierential equation. In

particular, we can add the transient solutions we found in 2.1 to the steady state solutions

we found in 3.1.1 to obtain solutions to the full time-dependence of an oscillator.

This brings us to our method for finding the complete solutions: add the steady state

solution, which has no free integration constants, to the general solution for the undriven

damped oscillator, which has two undetermined integration constants. The two undeter-

mined integration constants are then chosen so that they satisfy the initial conditions, i.e.

the initial position and velocity of the oscillator. You might worry that in doing this, weve

missed something, that there is some other term we may have left out. Well, you neednt

worry. There is a very important theorem concerning linear dierential equations that says

that if you have found a solution to a dierential equation that satisfies all the initial con-

ditions, then you have you have found the solution. Its called the Uniqueness Theorem. It

is discussed in nearly every introductory text on dierential equations.

As an example of how this works, consider the case where a sinusoidal forcing term

is applied starting at t = 0 to a oscillator that is initially at equilibrium: x(0) = 0 and

x(0) = 0. Lets consider the underdamped case where Q > 12 , so we will use the solutions

from 2.1.1. Adding together steady state and the transient general solutions, Eqs. (3.24)

and (2.22), respectively, with, as yet, undetermined integration constants A1 and A2 , gives

t

65 Forced oscillators and resonance

t/2

x(t) = AR cos !t + AI sin !t + e (A1 cos !u t + A2 sin !u t) , (3.44)

where !u is given by Eq. (2.15) and ! is the drive frequency. The amplitudes AR and AI

from the steady state solution can be written in terms of the real and imaginary parts of

the response function (!) using Eqs. (3.25) and (3.26): AR = F0 0 (!) and AI = F0 00 (!).

Applying the initial conditions that x(0) = 0 and x(0) = 0 to Eq. (3.44) yields

x(0) = F0 0 (!) + A1 = 0 (3.45)

00

x(0) = !F0 (!) A1 + !u A2 = 0 . (3.46)

2

Solving for A1 and A2 gives

A1 = F0 0 (!) (3.47)

!

A2 = F0 00 (!) F0 0 (!) . (3.48)

!u 2!u

Substituting A1 and A2 into Eq. (3.44) gives the solution for the specified initial conditions

( " !#

0 t/2

x(t) = F0 (!) cos !t e cos !u t + sin !u t

2!u

" #)

!

+ 00 (!) sin !t e t/2 sin !u t . (3.49)

!u

Figure 3.6 shows the solutions given by Eq. (3.49). The upper plot shows the results for

a driving frequency of ! = 0.9 !u . Initially, both the (nearly) resonant frequency !u and

the drive frequency ! are present in the response. Because they are nearly equal, the two

frequencies produce beats with a period of T B = 1/ where = |R | = |!u !|/2. As

time passes, the transient response at frequency !u dies out, owing to the e t/2 damping

terms in Eq. (3.49), and the beats fade away. The damping occurs on a time scale of 2/ or

equivalently after about Q = !0 / oscillations (here, Q = 64).

The lower plot in Fig. 3.6 shows the results for a driving frequency at resonance ! = !u .

In this case there are no beats because the oscillator is driven at the resonant frequency. In

this case, the steady state response builds up monotonically on a time scale of 2/ (or Q

oscillations) towards the steady state solution.

Before moving on to applications of resonators, we pause briefly to reflect on the fact

that only the transient solution depends on the initial conditions; the steady state solution

is completely independent of the initial conditions. This means that no matter what the

starting conditions are, if a system is driven at a particular frequency and forcing, it will

end up in the same state. This is a property of linear oscillators, one that can be dramatically

violated for nonlinear oscillators.

Thus far we have considered the case that the periodic forcing is applied directly to an

oscillating mass. In many systems the forcing is applied indirectly, for example, by causing

the point from which a spring or pendulum is suspended to oscillate back and forth or up

t

66 Resonance

z(t)

3

2

1

0 t/T

10 20 30 40 50 60 70 80

1

2

3

z(t)

15

10

0 t/T

10 20 30 40 50 60 70 80

5

10

15

t

Fig. 3.6 Transient response for an oscillator starting from rest (z = 0 and z = 0) and a quality

factor Q = 64. (a) The driving frequency is ! = 0.9 !u , slightly off resonance, and

beats appear at the difference frequency = |R | = |!u !|/2. (b) The driving

frequency is ! = !u , exactly at resonance. The response is much larger for the drive

frequency at resonance (note the difference in the z scale for the two plots). The time

axis in both cases is expressed in units of the natural period T = 2/!0 of the

oscillator.

and down. Resonance is still observed in such systems, but the asymptotic low and high

frequency behavior can be quite dierent from the case of direct forcing.

As an example of indirect forcing, lets consider a pendulum, illustrated in Fig. 3.7,

whose pivot point is forced to oscillate horizontally back and forth according to x p =

x0 cos !t. The pendulum is free to rotate about the pivot point. The mass at the end of

the pendulum moves a distance s = l along an arc like a normal pendulum. For small

amplitude oscillations the motion of the mass at the end of the pendulum is to a good

approximation parallel to the displacement of the pivot point, so that the absolute position

of the mass is given by x p + s. The damping force on the pendulum is taken to be friction

from the hinge at the pivot point, and is proportional to s, the velocity of the mass relative

to the pivot point (not the fixed reference frame). In this case, the equation of motion is

t

67 Forced oscillators and resonance

given by

d2 (x p + s) ds

m = mg (3.50)

dt2 dt

Rearranging terms and using s = l, we obtain,

1 d2 x p

+ + !20 = , (3.51)

l dt2

p

where !0 = g/l and = /m. Equation (3.51) bears a striking similarity to Eq. (3.4), the

equation of motion for the forced damped oscillator, but with the forcing term replaced by

x p /l, which is proportional to the acceleration of the pivot. Such a device is sometimes

used as a seismometer, a possibility we explore further in 3.1.6.

Equation (3.51) is valid for arbitrary accelerations x p , so long as remains sufficiently

small. We begin our analysis by taking the acceleration to be sinusoidal. Later, we will

consider sudden impulses. We write the sinusoidal displacement as the real part of a com-

plex exponential

i!t

x p = x0 cos !t = x0 Re{e }, (3.52)

and, as previously, will keep only the real part of any solution we obtain. Proceeding as

before, we assume a steady-state solution of the form (t) = Ae i!t , where, following the

same reasoning, the frequency of the steady state response is the same as that of the driving

force. Substituting our trial solution (t) = Ae i!t and Eq. (3.52) into Eq. (3.51) gives

x0

( !2 i ! + !20 )A = !2 , (3.53)

l

t

Fig. 3.7 Pendulum with a moving pivot point. Gray image shows the equilibrium position of

pendulum and the pivot point.

t

68 Resonance

i!t

where we have canceled the common factors of e . Solving for A, we obtain

(x0 /l) !2

A= . (3.54)

!20 !2 i !

Writing A = A0 ei , where a0 and are taken to be real, our solution takes the form

i(!t )

(t) = A0 e , (3.55)

where

(x0 /l) !2

A0 = q (3.56)

(!20 !2 )2 + !2 2

0 1

BB ! CCC

= tan 1 BB@ 2 CA . (3.57)

!0 !2

Figure 3.8 shows plots of Eqs. (3.56) and (3.57) for several values of Q = !0 / . The res-

onance curves look similar to those we found previously for the forced damped oscillator.

However, the asymptotic behaviors here are dierent: for ! !0 , the amplitude of the

oscillation goes to zero, i.e. A0 ! 0, while for ! !0 , the amplitude goes to a finite

value, i.e. A0 ! x0 /l. The behavior at low frequency is consistent with what you might

have anticipated, namely, if the pendulum is moved back and forth too slowly, the pendu-

lum doesnt oscillate but moves as a whole with its support: 0. When the pendulum

framework is moved back and forth very quickly at high frequencies, the mass on the end

of the pendulum basically remains (very nearly) still. Its angular amplitude A0 is just the

distance x0 the pivot point moves divided by the length of the pendulum. As a result, the

oscillations" in the reference frame of the pivot point are exactly out of phase with the

pivot point.

Our quantitative analysis of the problem has shown that too slowly" means at a fre-

quency much less than the natural oscillation frequency !0 of the pendulum. Of course,

we might have guessed this from the outset, as there are only two time scales in the prob-

lem, !0 1 and 1 , and it makes sense physically. If !0 , meaning the damping is

weak relative to the oscillations, the pendulum oscillates with respect to the pivot points.

If !0 , meaning the damping is strong relative to the oscillations, the pendulum just

moves with the pivot point.

Oscillators are frequently used as detectors that selectively filter out all signals except the

one of interest. One of the most common of all such detectors is the radio receiver, which

uses an electronic oscillator whose resonant frequency can be tuned to detect the one signal

that interests us among the many that are available. We discuss the radio receiver in 3.2.1.

t

69 Resonant filters and detectors

a0 (w )

15

Q = 16

10

Dw Q 1

0 w /w0

0 1 2 3

f (w )

p

Q = 16

p /2 Q = 1/2

w /w0

t

0

0 1 2 3

Fig. 3.8 Amplitude (top) and phase of pendulum response (t) to periodic shaking of the

pendulum for values of the quality factor Q from the darkest to the lightest curves: 12 ,

1, 2, 4, 8, 16. The dashed curve corresponds to Q = 1. The height of the peak of the

resonance curve (top) is given by A0 (!0 ) = Q( f0 /!20 ) and width ! is proportional to

1/Q in the limit that Q 1. Phase (bottom) of the response of the pendulum. The

curves corresponding to Q = 1 are shown as dashed lines.

Resonators as filters

Figure 3.9(a) shows that the amplitude of the response of the pendulum seismometer is

constantor flatabove the resonance frequency fR but falls o dramatically as the fre-

quency is reduced below fR . The pendulum, in addition to serving as a seismometer, also

filters the input signals, responding to those with frequencies greater than the resonant fre-

quency while suppressing signals with frequencies less than the resonant frequency. Such

a filter is called a high-pass filter because it passes through the filter only those signals at

frequencies higher than the resonant frequency.

By contrast, the driven damped pendulum that we considered in 3.1.1 acts as a low pass

filter. Its behavior can be seen in Fig. 3.2 and in Fig. 3.9(b) above where the amplitude of

the response is shown for smaller values of Q.

We can make another kind of filtera resonant filterby using either one of these two

t

70 Resonance

kinds of pendulums in a mode where we have Q 1. In this limit, the pendulums respond

only over a narrow range of frequencies centered around the resonant frequency.

All three kinds of filters, low-pass, high-pass, and resonant filters find wide applica-

tion in science and technology. In this section, we have studied devices that are in eect

mechanical filters. These same three kinds of filters are used in vastly dierent contexts,

however. There are electronic filters, optical filters, and acoustic filters. In fact, you can

find filters in almost any context in which you find oscillations. In the next section, we take

a look at radio receivers, which use resonant electronic filters. We will study the simplest

example, but one that is sufficient to illustrate the basic functioning filters in radios.

Radio signals permeate virtually every corner of the world we live in. At any given point in

space, the electric and magnetic fields associated with radio waves, coming from sources

near and far, fluctuate in time in a seemingly random manner, as illustrated in Figure

3.10(a). Some of the sources generate radio waves that, for lack of a better term, we call

noise. But many of those sources broadcast radio waves that, for lack of a better term,

we call information: music, news, and talk from commercial, nonprofit, and other sources.

The remarkable thing about radio and television receivers is that they are able to extract

from the morass that part of the signal that originates from a single source. A receiver does

this by filtering out essentially all the signals from sources not of interest. Figure 3.10(a)

shows what the time dependence of the electric field at a particular point in New York

City might look like; here we have included signals from more than 60 AM and FM ra-

dio stations. The modulation of the amplitude of the signal comes about because of the

interference between radio signals from the the various sources, each of which has its own

characteristic frequency. Because many of the broadcast frequencies are very near each

other, summing them produces the characteristic pattern of beats that are visible in the sig-

nal. Figure 3.10(b) shows that part of the same signal that includes only the signal coming

t

Fig. 3.9 (a) Amplitude vs. frequency for seismometer pendulum. The values of Q plotted

proceed Q = 2 (uppermost trace), 1, 12 (solid line), 14 , and 18 (lowermost trace). (b)

Amplitude vs. frequency for driven pendulum described in 3.1.1 and by Eq. (3.12).

t

71 Resonant filters and detectors

from a particular radio station, one that broadcasts at 99.5 MHz.2 It is the task of the radio

receiver to extract only that frequency from complex waveform shown in Figure 3.10(a).

The detection of radio signals is a complex process that depends on the size, shape, and

orientation of the antenna. We are not going to focus on that part of the detection of radio

waves. Suffice it to say that electric field of the the radio waves causes the electrons in the

antenna to oscillate back and forth producing a current whose time dependence follows the

electric field and thus resembles the signal shown in Figure 3.10(a). Our focus here is going

to be on how the electronics in a receiver selects the desired signal and filters out everything

else. The key, of course, is that each radio station broadcasts at a dierent frequency, so the

task of the receiver is to detect only the desired frequency and filter out the others. Such a

task is perfect for the electronic equivalent of the resonant filters introduced in the previous

sections.

At the heart of a radio receiver is a resonant filter, which in its simplest form consists of

an inductor, a capacitor, and a resistor connected to an alternating current or voltage source.

In the case of a radio receiver, it is the current, produced by radio waves, from the antenna

that acts as a current source. The current generated by radio waves is tiny in general and

thus must be amplified by a considerable amount to be processed and rendered audible.

In addition, a receiver must extract the audio signal from the radio wave, which it does

through a process called demodulation, which works dierently for AM and FM signals,

for example. Thus, a radio receiver must detect radio waves, selectively filter, amplify, and

2 99.5 MHz corresponds to a radio station in the FM or frequency modulated band. The audio signal is encoded

in small variations in the broadcast frequency. This means that the instantaneous broadcast frequency varies

by up to 20 kHz from the center broadcast frequency of 99.5 Mz, or about 0.02%.

t

Fig. 3.10 (a) Time dependence of the electric field of a radio signal at a particular point in

space. (b) Filtered radio signal. Note the difference in the scales of the electric field

E(t), and especially that the signal in part (b) is only a small fraction of the overall

signal shown in part (a).

t

72 Resonance

t

Fig. 3.11 Resonant circuit with a resistor, inductor, and capacitor in series.

demodulate (and amplify again!) in order for us to hear the music. These four processes

arent always entirely independent of each other but for the moment, we shall focus only

on the filtering of a radio signal such that the signal from only one radio station is detected.

As noted above, dierent radio stations broadcast at dierent frequencies. The commer-

cial AM band corresponds to radio frequencies between about 520 kHz and 1.6 MHz, with

dierent radio stations signals typically separated by at least 10 kHz. The commercial

FM band corresponds to radio frequencies between about 87 MHz and 108 MHz, with

dierent radio stations signals typically separated by at least 100 kHz. So we will want

to design an electrical resonant filter that can tune to these frequencies with a sufficiently

sharp resonance such that nearby radio signals are filtered out.

To begin, we consider the circuit shown in Fig. 3.11. The voltage source serves as our

signal," which, for the moment, we will take to have a sinusoidal time dependence:

V(t) = V0 cos !t . (3.58)

The voltage drop across the capacitor is simply the instantaneous charge q on the capacitor

divided by its capacitance C. Similarly, the voltage drop across the inductor is L dI/dt and

across the resistor is IR. Thus we have

q dI

V(t) = + L + IR = V0 cos !t. (3.59)

C dt

The current is related to the charge on the capacitor by I = dq/dt. Therefore, to obtain a

dierential equation we can readily analyze, we take the derivation of Eq. (3.59) and obtain

d2 I dI 1

L 2

+ R + I = !V0 sin !t . (3.60)

dt dt C

This has the same form as previous equations we have analyzed involving resonance, with

the current I as the dependent variable, L appearing in the place of mass, R in the place of

damping, and with C 1 providing the restoring term. Dividing through by L we obtain

d2 I R dI 1 !V0

+ + I= sin !t . (3.61)

dt2 L dt LC L

Comparing this to Eq. (3.2), rewritten below for convenience

d2 z dz F0

+ + !20 z = cos !t , (3.62)

dt2 dt m

1/2

we can identify the damping term as = R/L, the resonance frequency as !0 = (LC) ,

t

73 Resonant filters and detectors

ratio of two time constants, L/R and RC.

The steady state response of the current I must oscillate at the driving frequency, just

as we found previously for mechanical oscillations. In electronics, it is conventional to

assume that the current has the form I(t) = I0 ei!t instead of I(t) = I0 e i!t , so we follow

that convention here. Writing sin !t = iei!t and substituting into Eq. (3.60), we obtain

1

!2 L I0 ei!t + i!R I0 ei!t + I0 ei!t = !V0 iei!t . (3.63)

C

Dividing both sides by i!I0 ei!t and rearranging terms a bit yields

V0 1

= R + i!L + . (3.64)

I0 i!C

The terms on the right hand side have the units of resistance. Taken together, they constitute

what is commonly called the complex impedence of the circuit, which is often denoted as

Z. Factoring out R, we can write the complex impedance as

" #

V0 L 1

Z= = R 1 + i! + (3.65)

I0 R i!RC

" !#

! !0

= R 1 + iQ . (3.66)

!0 !

Thus, the current is given by

V0 i!t

I(t) = I0 ei!t = e (3.67)

Z

V0

= ei!t . (3.68)

R + i!L + 1/i!C

We can write I0 = V0 /Z in polar form I0 = |I0 |ei , where

2 !2 3 1/2

|V0 | 6666 2 ! !0 7777

|I0 | = 61 + Q 7 (3.69)

R 4 !0 ! 5

" !#

!0 !

= tan 1 Q . (3.70)

! !0

Note that the current goes to zero at zero frequency and to a finite value for ! !0 . Thus,

for the overdamped case when Q 1/2, the circuit would act as a high-pass current filter.

t

74 Resonance

Problems

3.1 A block of mass m = 100 g is hung from a spring whose force constant is k = 60 N/m.

The damping force on the block is given by Fd = 3, where = 3.0 N-m/s. The

mass is subject to an oscillatory driving force of F = F0 cos !t, where F0 = 1.5 N.

(a) Write down the dierential equation describing the steady state motion of the

block and solve for its amplitude and phase as a function of angular frequency.

Express your answers entirely in terms of symbols (no numbers).

(b) What are the resonant frequency and Q of the system?

(c) What is the amplitude of oscillation at the resonant frequency?

(d) At what frequencies is the amplitude of oscillation half of its maximum value?

(e) What are the limiting amplitudes of oscillation for ! much less than and much

greater than resonant frequency?

(f) What is phase dierence between the driving force and the displacement at res-

onance? for ! = !0 /2? for ! = 2!0 ? Give your answers in both radians and

degrees and find the time delay between the peak in the driving force and the

peak in the mass displacement.

3.2 A 0.30-kg mass hangs at equilibrium from a spring suspended from a rod. Pulling the

mass down 4.0 cm from its equilibrium position requires a force of 5.0 N to stretch

the spring.

(a) Letting go of the mass after displacing it 4.0 cm from its equilibrium position,

you observe that its amplitude decays to 0.50 cm after 10 oscillations. What is

the Q of this system?

(b) Suppose that the system is now driven with a oscillating driving force of F(t) =

F0 sin !t starting at time t = 0, where F0 = 0.50 N and the driving frequency

f = 2.5 Hz. The initial displacement of the mass is, once again, 4 cm below its

equilibrium position.

(i) Find an equation for the subsequent motion of the mass and then plot it as

a function of time.

(ii) What is the steady state amplitude the oscillations?

(iii) What is the initial period of the oscillations?

(iv) What is the ultimate period of the oscillations?

(v) How much energy is stored in the oscillating system at steady state?

(vi) How much energy is dissipated per period by the system at steady state?

3.3 Imagine an experiment similar to Millikans famous oil drop experiment in which

an oil droplet of mass m is suspended between two flat metal plates a distance h

apart and connected to a power supply that produces an oscillating voltage V(t) =

V0 cos !t across the plates. Further assume that that the oil droplet carries a charge of

q, such that it experiences an oscillating force F = (qV0 /h) cos !t from the oscillating

electric field. The damping force from the droplet moving through air is given by

t

75 Problems

Fd = 3, where is the damping constant and 3 is the droplet velocity. You may

ignore gravity in this problem. In this problem there is no spring-like restoring force.

(a) Write down the equation of motion for this system, including the oscillatory

driving force and proceed to solve it by assuming a trial solution of the form y =

Ae i!t . Solving the resulting equation, show that the steady-state displacement

of the droplet is given by x(t) = A0 (!) cos(!t ) where

qV0

A0 = p ,

m h ! !2 + 2

1

= tan ,

!

and = /m.

(b) The damping constant can be estimated from the equation 3d, where

1.8 10 5 Pa-s is the viscosity of air (in SI units) and d is the diameter

of the droplet. Assuming a droplet diameter of d = 1 mm, a plate separation of

h = 1 cm, and an oil mass density of 0.8 g/cm3 , find the amplitude of oscillation

for a drive frequency of ! that is equal to the damping rate /m and a voltage

amplitude of V0 = 100 V. Comment on the feasibility of such an experiment

given the numbers you find.

3.4 Consider an spring-and-mass oscillator such as that of Fig. 3.1, the top of whose

spring is moved up and down with an amplitude s(t) = a cos !t. Set up the equation

of motion for this system, and show that the steady-state solution for the displace-

ment x(t) for the weight at the end of the spring is

!20

x(t) = a q cos(!t ).

(!2 !20 )2 + !2 2

1 1/Q

Pd = F2 .

2m!0 0 !

! 2 0

+ 1

!0 ! Q2

Use this result to show that Pd [Eq. (3.33)] has its maximum at ! = !0 and that the

full width at half maximum is !fwhm = !0 /Q.

3.6 Consider a driven damped oscillator consisting of a mass hanging from a spring like

the one depicted in Fig. 3.1. Initially the system hangs at its equilibrium position.

At time t = 0 a sinusoidal oscillating force is applied to the mass: F(t) = F0 sin !t.

Following the method used in 3.1.5, find the subsequent motion of of the oscillator

for a system with Q = 64 and a driving frequency of ! = 0.9!u . Take the period of

the oscillator to be T = 2/!0 = 1 s. Compare your result with the result obtained

in 3.1.5 for a cosine driving force at times t QT and for times t QT .

3.7 Repeat Problem 3.6 except instead of applying a sinusoidal force starting at t = 0,

solve the problem for the case where the point from which the (massless) spring

hangs starts oscillating up and down at t = 0 with a time-dependent vertical dis-

placement given by y p (t) = y0 sin !t. Take the damping force to be Fd = y. Take

t

76 Resonance

the period of the oscillator to be T = 2/!0 = 1 s. Compare your result with the

result obtained in Problem 3.6 at times t QT and for times t QT .

3.8 The seismometer described in 3.1.6 can detect seismic motion in the horizontal

direction. To detect motion in the vertical direction, we need an oscillator capable of

motion in the vertical direction. A simple mass hanging from a suitable spring would

work, assuming the system is constrained to move only in the vertical direction.

(a) Write down Newtons second law for this system and obtain the equation of

motion assuming the Earth and oscillator support oscillates vertically according

to the equation yE (t) = a cos !t, where a is the amplitude of the motion. You

should follow the same kind of procedure used in 3.1.6 to set up the swing-gate

seismometer by defining coordinates relative to a fixed reference frame (not the

Earth!).

(b) Find the steady state motion, including the frequency-dependent amplitude A0 (!)

and phase (!) relative to the driving force. Plot A0 (!) and (!) as a function

of frequency for Q = 14 , 12 , 1, 2, 4, 8 and 16.

(c) Find and plot the response of the seismometer to a sudden impulse at time t = 0

assuming the mass is initially at rest at its equilibrium position. Choose values

for the spring constant, damping constant, and mass so that the instrument is

sensitive to frequencies from 0.1 to 1 Hz.

(d) Does this seismometer function as a low pass or high pass filter? Briefly explain.

3.9 In this problem, we consider an electrical resonant circuit consisting of two resistors,

a capacitor, and an inductor, driven by a sinusoidal voltage source V(t) = V0 cos !t

across the open terminals shown in the figure below.

t

Fig. 3.12 Resonant RLC circuit.

(a) Find the dierential equations describing for the current I in each branch of the

circuit and then solve for the voltage across the capacitor as a function of time.

Hint: Assume, as we did in 3.2.1, that the current has a complex amplitude and

is proportional to ei!t in each branch of the circuit and apply Kirchhos law for

the current at each junction.

(b) What are the resonant frequency and Q for this circuit in terms of R, L, and C?

(c) On a single graph, plot the amplitude and phase of the voltage across the capaci-

tor in the steady state as a function of !/!0 assuming a values of Q = 1, 10, and

100. Be sure to label which curves belong to which values of Q.

t

77 Problems

3.10 (French 4-3) An object of mass 0.2 kg is hung from a spring whose spring constant

is 80 N/m. The body is subjected to a resistive force given by b3, where 3 is the

velocity (m/sec) and b = 4 N-m 1 s.

1 Set up the dierential equation for free oscillations of the system and find the

period of such oscillations [Ans: T = 5 p 3 = 0.36s]

2 The object is subjected to a sinusoidal force given by F(t) = F0 sin !t, where

F0 = 2 N and ! = 30 s 1 . In the steady state, what is the amplitude of forced

oscillations? [Ans: A = 1.28 cm]

3.11 (French 4-11) Consider a damped oscillator with m = 0.2 kg, b = 4 N-m 1 -s and

k = 80 N/m. Suppose the oscillator is driven by a force F = F0 cos !t, where F0 = 2

N and ! = 30 s 1 .

1 What are the values of A and of the steady-state response described by x(t) =

A cos(!t ) [Ans: A = 1.3 cm, = 130 ]

2 What is the mean power input. [Ans: P = 18/61 = 0.295 J/s.]

3 How much energy is dissipated against the resistive force in one cycle? [Ans:

Energy/cycle = PT = 0.062 J]

4 Normal Modes

Up to this point, we have considered oscillations of systems with with only a single degree

of freedomthe angle a pendulum makes with the vertical, the distance between two

atoms in a diatomic molecule, etc. However, most things that oscillate have many degrees

of freedom and typically oscillate at many dierent frequencies. Consider a string on a

guitar or piano. It can vibrate at its fundamental frequency or its first, second, or higher

overtones. The same is true of every other kind of musical instrument, or in fact, almost

any other object that vibratesa plate, a goblet, a bridge, or even a skyscraper. Typically,

for each frequency of vibration, there is a characteristic pattern of displacements of the

vibrating object. These patterns of displacements and the characteristic frequency that goes

with each of them are known as normal modes.

In this chapter we begin our study of normal modes with examples of systems that have

only two (or three) degrees of freedom, as these serve to introduce the basic concepts. Once

you understand those systems, we will move on to systems with more degrees of freedom,

where the real power of normal mode analysis, both conceptually and mathematically,

is revealed. As might be expected, when the number of degrees of freedom increases,

the algebra can become very tedious. To keep the algebra manageable, we introduce the

the tools of linear algebra: matrices, vectors, and eigenvalues. The mathematics of linear

algebra also introduces a conceptual framework for better understanding the physics of

systems with many degrees of freedom. Dont worry if you are not too familiar with linear

algebra. We will introduce the tools and concepts of linear algebra as we need them.

We begin our exploration of normal modes by considering two identical pendulums con-

nected by a spring, as shown in Fig. 4.1. The spring connecting the two masses is un-

stretched when the two pendulums hang vertically, that is, when 1 = 2 = 0. Gravity

provides a restoring force for the pendulums of F1g = mg sin 1 and F2g = mg sin 2 ,

respectively. When stretched or compressed, the spring exerts a force F1s ' k(x2 x1 ) on

mass 1 and F2s ' k(x2 x1 ) on mass 2 in the limit that 1, where we can neglect the

eects of the vertical extension of the spring. The sign convention is chosen so that when

78

t

79 Systems with two degrees of freedom

1 2

t

x1 x2

Fig. 4.1 Two identical pendulums of length l and mass m connected by a spring with spring

constant k.

the spring is stretched, it exerts a force to the right on mass 1 and to the left on mass 2.

Summing the forces on each mass, we obtain the following equations of motion:

x1

F1 = mg + k(x2 x1 ) = m x1 (4.1)

l

x2

F2 = mg k(x2 x1 ) = m x2 , (4.2)

l

where we have used the approximations sin 1 ' x1 /l and sin 2 ' x2 /l, which are valid

when x1 , x2 l. Collecting terms together and rearranging gives

mg

m x1 = + k x1 + kx2 (4.3)

l

mg

m x2 = kx1 + k x2 , (4.4)

l

Equations (4.3) and (4.4) are the equations of motion for two pendulums connected by

a spring. They are coupled equations in that the dynamical variables x1 and x2 appear in

both equations. Physically, this means that exciting pendulum 1, say by giving it some

finite displacement, x1 , 0, or finite velocity, 31 , 0, will in general cause pendulum 2 to

move so that x2 , 0 and 32 , 0. Thus we say that the dynamical variables x1 and x2 are

coupled. Mathematically, the fact that these equations are coupled means that we will need

to solve for both x1 and x2 at the same time.

We are going to solve Eqs. (4.3) and (4.4) in two stages. First we will find the natural

frequencies of oscillation, the so-called normal frequencies. This particular system has two

degrees of freedom, the coordinates x1 and x2 of the two pendulums. Because there are two

degrees of freedom, there will be two normal frequencies. In general, the number of normal

frequencies is equal to the number of degrees of freedom, although in some cases there can

be fewer normal frequencies.

Second, we will find equations of motion for x1 and x2 . We will do this using a some-

what indirect approach that benefits from the hindsight that comes from others having

solved these kinds of problems many times before. The approach consists of defining new

dynamical variables q and q that are linear combinations of x1 and x2 . By choosing those

linear combinations wisely, we will be able to turn the two coupled dierential equations,

t

80 Normal Modes

Eqs. (4.3) and (4.4), into two independent uncoupled equations that are in the precise form

of the equation of motion for a single simple harmonic oscillator that we encountered in

Chapter 1 (see Eq. (1.5)). And we will be able to use the same sinusoidal solutions that we

developed in Chapter 1.

Normal frequencies

To solve Eqs. (4.3) and (4.4), we follow the same procedure we used for the case of a single

oscillator and assume complex sinusoidal solutions of the form

i!t

x1 = a1 e (4.5)

i!t

x2 = a2 e . (4.6)

Substituting Eqs. (4.5) and (4.6) into Eqs. (4.1) and (4.2) and collecting terms gives

!

g k k

+ !2 a1 a2 = 0 (4.7)

l m m

!

k g k

a1 + + !2 a2 = 0 . (4.8)

m l m

These equations are satisfied only when the determinant of the matrix formed by the coef-

ficients of the amplitudes a1 and a2 is zero,1 that is, when

=0 (4.9)

k/m g/l + k/m !2

which gives

!2 !2

g k k

+ !2 =0. (4.10)

l m m

Equation (4.10) is known as the secular or characteristic equation for the system of cou-

pled equations formed by Eqs. (4.7) and (4.8). Solving the secular equation2 for ! gives

q

! = (!2p + !2s ) !2s , (4.11)

p p

where ! p = g/l and ! s = k/m. Equation (4.11) yields two frequencies, which we write

out explicitly

r

g

! = ! p = (4.12)

l

r

2 2 1/2 g 2k

! = (! p + 2! s ) = + . (4.13)

l m

1 You are asked to prove this assertion, which is a special case of Cramers rule, in Exercise 4.1.1. Its handy to

remember it! For an introduction and review of matrices and linear algebra, see Appendix A.

2 You can blindly expand Eq. (4.10) and solve it using the quadratic formula, or you can save yourself a lot of

work by noticing that it consists of two perfect squares, allowing you to simply write down the solution by

inspection, saving work and avoiding algebra errors.

t

81 Systems with two degrees of freedom

These two frequencies ! and ! are the normal frequencies of this system. To under-

stand their physical significance, we need to find the normal coordinates q and q and the

normal modes of the system.

Exercise 4.1.1 Cramers Rule: Given a pair of coupled linear equations of the form

B21 a1 + B22 a2 = 0 , (4.15)

!

B11 B12

B= , (4.17)

B21 B22

then the determinant of the 2 2 matrix B is written as |B| or det B and is defined as

B11 B22 B12 B21 . Thus, Eq. (4.16) can be written as |B | = det B = 0.

Normal coordinates

An interesting and extremely useful feature of systems of coupled linear oscillators is that

it is always possible to define new coordinates q and q that are linear combinations of the

original coordinates x1 and x2 such that the equations of motion for these new dynamical

variables are uncoupled. Analyzing the composite system in terms of these new dynamical

variables results in an enormous mathematical and conceptual simplification of the system.

The new dynamical variables are sufficiently important that they are given a special name:

normal coordinates. Fancy words! Lets illustrate what they mean by finding the normal

coordinates of the system of coupled pendulums we have been studying.

The idea is to recast Eqs. (4.1) and (4.2) so that we obtain two equations, each of which

involves a single dynamical variable that does not appear in the other equation. This is

readily accomplished by summing and subtracting Eqs. (4.1) and (4.2). Summing Eqs.

(4.1) and (4.2) gives

x1 + x2

mg = m( x1 + x2 ) . (4.18)

l

Subtracting Eq. (4.1) from Eq. (4.2) gives

x2 x1

mg 2k(x2 x1 ) = m( x2 x1 ) , (4.19)

l

Collecting and rearranging terms, these equations become

g

( x1 + x2 ) + (x1 + x2 ) = 0 (4.20)

!l

g 2k

( x2 x1 ) + + (x2 x1 ) = 0 . (4.21)

l m

t

82 Normal Modes

q = x1 + x2 (4.22)

q = x2 x1 , (4.23)

g

q + q = 0 (4.24)

l!

g 2k

q + + q =0. (4.25)

l m

The new coordinates q and q are linear combinations of the original coordinates x1 and

x2 . Moreover, with this choice of q and q , the equations of motion, Eqs. (4.22) and (4.23),

become decoupled: Eq. (4.22) is the equation of motion for q only, and Eq. (4.23) is the

equation of motion for q only. This is important. It means that the dynamics of q , that

is how q develops in time, does not depend on q , and vice versa; q (t) and q (t) are

completely independent of each other. The new dynamical variables q and q are called

the normal coordinates of the system.

The equations of motion Eqs. (4.22) and (4.23) for the normal coordinates q and q are

the same as the equation of motion for a simple harmonic oscillator (cf. Eq. (1.5)),

q + !2 q = 0 (4.26)

2

q + ! q = 0 , (4.27)

where ! and ! are constantsthe oscillation frequenciesof the dierent normal modes.

This is a general feature of normal coordinates: their corresponding equations of motion are

always the equation of motion of the simple harmonic oscillator. The only distinguishing

feature is that each normal coordinate has its own characteristic frequency of oscillation.

Comparing Eqs. (4.22) and (4.23) with Eq. (1.5), we can immediately read o the oscil-

lation frequency corresponding to each normal coordinate:

p q oscillates at the frequency

! = sqrtg/l; q oscillates at the frequency ! = g/l + 2k/m. Of course, these are just

the frequencies of the normal modes we identified in Eqs. (4.12) and (4.13). This is also

a general feature of normal coordinates: the natural oscillation frequency of each normal

coordinate corresponds to one of the normal frequencies of the system.

The two frequencies ! and ! and the two normal coordinates q and q correspond to

two separate modes of oscillation of the system of two coupled pendulums. In the first

mode, both pendulums oscillate at a frequency

p !p ; in the second mode, both pendulums

oscillate at a frequency ! . Since ! = g/l < g/l + 2k/m = ! , the mode is slower

than the oscillation mode. Lets examine each of these modes individually.

t

83 Systems with two degrees of freedom

We start by examining the slow mode, which oscillates at the frequency ! . The solutions

are given by summing the positive and negative frequency ! solutions given by Eqs.

(4.5) and (4.6)

i! t

x1 (t) = a1 e + a1 ei! t (4.28)

i! t

x2 (t) = a2 e + a2 ei! t , (4.29)

where we recall that the two terms in each sum must be complex conjugates of each p other

for the solutions x1 (t) and x2 (t) to be real (see 1.5.2). Substituting ! = ! = g/l into

Eq. (4.7), we find that a1 = a2 for the (slow) mode. With a1 = a2 , Eqs. (4.28) and (4.29)

are identical. That is, for the slow mode, x1 (t) = x2 (t).

The equation of motion for the slow mode is given by Eq. (4.24) or equivalently Eq.

(4.26). As the equation of motion is identical to that of a single simple harmonic oscillator,

so the general solution can be written as

where we chosen to insert a factor of 2 into the definition of the amplitude A (see Eq.

(1.71) and 1.5.2). Since q (t) = x1 (t) + x2 (t) and x1 (t) = x2 (t), the general solutions for

the equations of motion for the individual masses oscillating in the slow mode are

According to Eq. (4.31), the two masses have exactly the same equation of motion, mean-

ing that they oscillate in phase with each other at a frequency ! . Because they move in

phase with each other, the spring connecting them is neither stretched nor compressed; the

distance between the masses remains

p unchanged. Thus, gravity provides the only restoring

force and the frequency ! = g/l of the coupled system is the same as for the uncoupled

system.

t

x1 x2 x1 x2

t

84 Normal Modes

Following the same procedure we used for analyzing the slow mode, the solutions for the

fast mode are given by summing the positive and negative frequency ! solutions given

by Eqs. (4.5) and (4.6)

i! t

x1 (t) = a1 e + a1 ei! t

(4.32)

i! t

x2 (t) = a2 e + a2 ei! t . (4.33)

p

Substituting ! = g/l + 2k/m into Eq. (4.7), we find that a1 = a2 for the (fast) mode.

With a1 = a2 , x1 (t) = x2 (t).

The equation of motion for the fast mode is given by Eq. (4.25), or equivalently by Eq.

(4.27), and has the same general solution

q (t) = 2B cos(! t ). (4.34)

Since q (t) = x1 (t) + x2 (t) and x1 (t) = x2 (t), the general solutions for the equations of

motion for the individual masses oscillating in the slow mode are

x1 (t) = x2 (t) = B cos(! t ). (4.35)

For the fast mode, we see that the two pendulums oscillate synchronously exactly out of

phase with each other, as shown in the right panel in Fig. 4.2. In the fast mode, the spring

is alternately stretched and compressed each cycle. In this case, both the spring and gravity

provide the restoring force, which is reflected in the equation for the normal frequency

! = (g/l + 2k/m)1/2 = (!2p + 2!2s )1/2 .

Because the equations of motion are linear, linear combinations of the solutions, including

those of dierent normal modes, are also solutions to the equations of motion. This is

the principle of linear superposition. In fact, all possible solutions of any coupled linear

oscillator can be expressed as a linear combination of the solutions to the individual normal

modes.

This brings us to the practical question of how best to find the solution to a coupled

oscillator problem for a particular set of initial conditions. In general, the simplest way to

solve such a problem is (i) first find the normal coordinates of the system, (ii) then express

the initial conditions in terms of those normal coordinates, (iii) and finally, find the solution

for each normal mode subject to the initial conditions for that normal coordinate. Once you

have the solutions in terms of the normal coordinates q (t) and q (t), find the equations for

x1 (t) and x2 (t) by inverting the equations for q (t) and q (t).

Lets illustrate this procedure with an example. Before doing so, however, we note that

each normal mode is governed by an equation of motion for a simple harmonic oscillator:

q + !2 q = 0. (4.36)

Because the equation of motion for each normal coordinate is the equation of motion of a

simple harmonic oscillator, the solution for each normal mode can be written in any one of

t

85 Systems with two degrees of freedom

several forms. First there is the complex exponential form which is useful for setting up the

problem and for finding the normal frequencies and normal modes. But for constructing

useful solutions, it is usually better to use real functions. The two most common forms that

we have employed are

q (t) = a1 cos ! t + a2 sin ! t (4.37)

q (t) = A cos(! t ), (4.38)

where refers to a particular normal mode. In either case, there are two integration con-

stants, consistent with the second order equation of motion Eq. (4.26), to determine from

the initial conditions, the amplitudes a1 and a2 in the first case and the amplitude A and

phase in the second. Which form is used to solve a particular problem is to some extent

a matter of taste as the two descriptions are completely equivalent. However, at the final

stage of finding a particular solution to a normal mode problem with specific initial condi-

tions, it is generally simpler to use sines and cosines rather than the cosine with a phase.

We illustrate how this works by applying it to obtain solutions for two coupled pendulums.

Example 4.1 Find the equations of motion of two pendulums of mass m coupled by

a spring with a spring constant k for the initial conditions x1 (0) = 1, x2 (0) = 0, and

v1 (0) = v2 (0) = 0.

Solution

We solve this problem in terms of the normal coordinates, which we define using the un-

normalized form given by Eqs. (4.22) and (4.23), and using the sine and cosine form of the

solutions given by Eq. (4.38):

q (t) = A1 cos ! t + A2 sin ! t (4.39)

q (t) = B1 cos ! t + B2 sin ! t . (4.40)

From Eqs. (4.22) and (4.23), the initial conditions for the normal coordinates are

q (0) = A1 = x1 (0) + x2 (0) = 1 (4.41)

q (0) = B1 = x2 (0) x1 (0) = 1 . (4.42)

Similarly, the initial conditions for the velocities gives

q (0) = ! A2 = x1 (0) + x2 (0) = 0 (4.43)

q (0) = ! B2 = x2 (0) x1 (0) = 0 , (4.44)

which means that A2 = B2 = 0. Therefore, the full solutions for the normal modes are:

q (t) = cos ! t (4.45)

q (t) = cos ! t . (4.46)

Using Eqs. (4.22) and (4.23), these solutions give

x1 (t) = 12 [q (t) + q (t)] = 12 [cos ! t + cos ! t] (4.47)

1 1

x2 (t) = 2 [q (t) q (t)] = 2 [cos ! t cos ! t] . (4.48)

t

86 Normal Modes

1.0

0.5

0.0 t

x1

5 10 15 20 25 30 35 40

0.5

1.0

1.0

0.5

0.0 t

x2

5 10 15 20 25 30 35 40

t

0.5

1.0

Fig. 4.3 Displacements of two pendulums coupled by a spring for the case where the natural

periods of the pendulums and spring are T p 2/! p = 3 s and T s 2/! s = 3 s,

respectively, which, from Eqs. (4.12) and (4.13), corresponds to T = 2/! = 3 s and

T = 2/!p = 1.73 s. For thep numerical values chosen here,

T /T = 1 + (T p /T s )2 = 2. Because this ratio is an irrational number, the

displacements of the two pendulums never repeat.

These solutions can produce quite complicated trajectories for the two pendulums, as

shown, for example, in Fig. 4.3. For this case, we have chosen the natural frequencies ! p

and ! s (or periods T p and T s ) of the spring and pendulums to be comparable to each other.

In this limit the force exerted by the spring on the two mass is comparable to the force of

gravity, and we say that the two pendulums are strongly coupled. In the next section, we

examine the particularly interesting limit of weak coupling.

It often turns out that the dierent degrees of freedom of a system are only weakly coupled

to each other. By this we mean that the characteristic energy (or force) with which dierent

degrees of freedom interact is small compared to the characteristic energies (or forces)

associated with the oscillations of those degrees of freedom. This arises, for example, for

the case of the two pendulums coupled by a spring when the characteristic energy E s of

compressing the spring is much smaller than the gravitational energy E p associated with

the pendulums rise. In the limit of small amplitude oscillations, these two energies can be

written as E s m!2s a2 and E p m!2p a2 (see Eq. (1.53)). The weak coupling limit thus

corresponds to cases for which E s E p , or equivalently, !2s /!2p 1.

The trajectories for two weakly coupled pendulums are shown in Fig. 4.4. The trajecto-

ries resemble the familiar pattern of beats observed when two sound waves of very nearly

t

87 Systems with two degrees of freedom

1.0

0.5

0.0 t

x1

0 5 10 15 20 25 30 35 40

0.5

1.0

1.0

0.5

0.0 t

x2

0 5 10 15 20 25 30 35 40

t

0.5

1.0

by a spring. Gray

dashed

lines

show modulated amplitudes: top, cos 12 ! t and bottom, sin 12 ! t .

the same frequencies interfere. Similarly, a beat pattern appears in the trajectories of cou-

pled two pendulums when the normal mode frequencies are very nearly the same. In the

present case, the dierence in the normal mode frequencies is given by

!=! ! = (!2p + 2!2s )1/2 !p (4.49)

1/2

= !p 1 + 2!2s /!2p 1 (4.50)

h i

= ! p 1 + !2s /!2p 1 4 4

2 ! s /! p + ... 1 (4.51)

!

!s

' !s (4.52)

!p

where we have used the binomial expansion Eq. (1.90) in the penultimate step and have

dropped the higher order terms in the expansion in the final step, which is a good approx-

imation in the limit that !2s /!2p 1. In that limit, it is also clear that ! ! , ! , and

therefore that the two normal mode frequencies ! and ! are nearly the same. Defining

! = 12 (! + ! ), we can rewrite Eqs. (4.47) and (4.48) as

h i h i

x1 (t) = 12 cos (! 12 !)t + 12 cos (! + 12 !)t (4.53)

1

= cos 2 ! t cos !t (4.54)

h i h i

x2 (t) = 12 cos (! 12 !)t 1

2 cos (! + 12 !)t (4.55)

1

= sin 2 ! t sin !t

, (4.56)

where we have used the cosine angle addition formulas to obtain Eqs. (4.54) and (4.56).

Equations (4.54) and (4.56) show that x1 (t) and x2 (t) oscillate at the mean frequency !

of

the two normal modes with a sinusoidally modulated amplitude with nodes spaced by a

time 2/ !.

t

88 Normal Modes

The trajectories plotted in Fig. 4.4 show an interesting pattern where the two oscillators

take turns oscillating; when one oscillates with its maximum amplitude, the other is

essentially still. This is one of the hallmarks of two weakly coupled oscillators and can be

observed in a very wide variety of systems.

From the previous discussion it should be clear that energy is passed back and forth be-

tween mass 1, mass 2, and the spring. Consider the case of the weakly coupled pendulums

discussed in the previous section. The gravitational potential energy and kinetic energy of

mass 2 are zero at t = 0 and again at t = 2/ !, since the displacement and velocity of

mass 2 are zero at those times. On the other hand, the gravitational potential energy of

mass 1 is at a maximum, as it is maximally displaced from its equilibrium position at those

times. The potential energy of the spring is non-zero as it is stretched from its equilibrium

position. As time passes, all these energies change, even as the total energy of the system

must remain constant, since no energy is dissipated.

We can make an even stronger statement about energy conservation: the total energy

associated with each normal mode is constant in time. This follows from the fact that

normal modes are independent of each other. Because they are independent of each other,

energy is not passed back and forth between the modes but is independently conserved for

each mode. The amount of total energy stored in each normal mode depends on the initial

conditions, which determine the degree to which mode is excited.

Determining exactly how much energy is to be associated with each mode can be tricky.

For example, what is the kinetic energy of the mode as a function of time? For a single

particle of mass m and velocity 3, the kinetic energy is simply 12 m32 . For the mode, we

might take the velocity to be equal to q (t), but what mass should we use? The mass m of a

single particle? or the mass 2m of two particles (since two masses are involved)? or some

other value? And what is the potential energy for mode 1? It seems clear that we should not

include the potential energy of the spring, since the mode does not involve stretching the

spring. But how do we divide the gravitational potential energy of the two masses divided

between the and modes?

Lets start by determining the kinetic energy K (t) as a function of time for the mode.

We can do this by calculating the kinetic energies of particles 1 and 2 assuming that only

the mode is excited initially, since what happens in the mode is independent of what

happens in the mode. For the spring-coupled pendulums, the displacements of the two

masses are given by Eqs. (4.47) and (4.48)

1

) x1 (t) = 2 q (t) (4.58)

1

x2 (t) = 2 [q (t) q (t)] (4.59)

1

) x2 (t) = 2 q (t) . (4.60)

The kinetic energy in the mode is just the sum of the kinetic energies of masses 1 and 2

t

89 Matrix formulation of normal modes

under the assumption that the mode is not excited. Thus, we obtain

2 2

= 12 m 12 q + 12 m 12 q (4.62)

1 2

= 2 (m/2)q . (4.63)

Based on our choice of normal coordinates, the eective mass me associated with the q

mode is m/2 according to Eq. (4.63).

The gravitational potential energy in the mode is obtained by summing the potential

energy of particles 1 and 2 assuming the mode is not excited:

= mgl(1 cos 1 ) + mgl(1 cos 2 ) (4.65)

h i h i

= mgl 1 1 12 12 + ... + mgl 1 1 1 2

2 2 + ... (4.66)

' 12 mgl 12 + 22 (4.67)

mg 2

' x1 + x22 (4.68)

2l

mg 1 2 1 2

' q + 2 q (4.69)

2l 2

1 mg 2

= q , (4.70)

2 2l

where we have used the small-amplitude approximation. Equation (4.70) is in the usual

form 12 ke x2 for the potential energy of a simple harmonic oscillator with an eective

spring constant ke = (m/2)(g/l) = me !2 , which is consistent with the usual relation for

an ideal spring simple

p harmonic oscillator with an eective mass of m/2 and an oscillation

frequency of ! = g/l. We note that our results for the eective mass me and eective

spring constant ke depend on our choice of normal coordinates, which are determined

only to within an overall multiplicative constant.

Exercise 4.1.2 Following the same procedures used above, find the small angle expres-

sions for K (t) and U (t), the time-dependent kinetic and potential energies of the

mode for the spring-coupled pendulums. In particular, show that the eective mass of

the mode is given by m/2, and that the expression for the eective spring constant

is consistent with this value and the value of ! . Then show by explicit calculation

that the total energy for each of the modes E = K (t) + U (t) and E = K (t) + U (t)

are constants independent of time.

The coupled pendulum problem discussed in the previous sections is actually one of the

simplest examples of coupled oscillators. Because it is so simple, it is relatively easy to

t

90 Normal Modes

analyze. The normal coordinates q and q turn out to be about the simplest set of linear

combinations of the particle coordinates x1 and x2 imaginable: q = x1 +x2 and q = x1 x2 .

So it was fairly easy to guess their form just by looking at the equations of motion, Eqs.

(4.3) and (4.4).

Unfortunately, finding the normal modes for most coupled oscillator problems is not so

simple. Moreover, the algebra associated with normal modes can get fairly tedious and

complicated. While tedious algebra is simply a fact of life sometimes, it can obscure the

physics and impede conceptual development of a subject. Introducing more compact math-

ematical notation, properly conceived, can improve matters. Consider, for example, the in-

troduction of vectors into physics. Using vector notation, we can write Newtons second

law as

F = ma , (4.71)

which with one equation neatly summarizes three equations written in algebraic form:

Writing Newtons second law in vector form can greatly simplify algebraic manipulations,

even if eventually we might need to use the component forms to do some of the mathemat-

ical manipulations.

There is another perhaps less apparent advantage to writing Newtons second law in

vector form. It helps us to think about the physics of Newtons second law in terms of

vectorsarrowsthat point some direction in space without reference to any particular

coordinate system. Thus we dont have to get bogged down in defining a coordinate system

and thinking only in components along one axis or another. We can think about the essential

elements of the problem geometrically, using helpful and compelling pictures. Enabling

us to think on this conceptual level is every bit as important as facilitating mathematical

manipulations.

For these reasons, both practical and conceptual, we embark below on a reformulation

of normal modes in terms of matrices and the mathematics of linear algebra.

Consider two equal masses m constrained to move along a frictionless surface as shown in

Fig. 4.5. The equations of motion for the two masses are

m x1 = ka x1 + kb (x2 x1 ) (4.73)

m x2 = kb (x2 x1 ) . (4.74)

t

Fig. 4.5 Horizontal masses on a frictionless surface coupled by two springs.

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91 Matrix formulation of normal modes

To determine the sign on the (x2 x1 ) terms, note that when (x2 x1 ) > 0, the kb spring

is stretched and it pushes mass 1 to the left in the negative direction and mass 2 to the

right in the positive direction. Collecting the x1 and x2 terms together, we rewrite these two

equations in standard form

m x1 = (ka + kb )x1 + kb x2 (4.75)

m x2 = kb x1 kb x2 . (4.76)

This equation can be written in compact matrix form with the introduction of the stiness

matrix k and the column vector x

! !

ka + kb kb x1

k= , x= . (4.77)

kb kb x2

With these definitions, the two equations of motion can be expressed as a single matrix

equation

1

x(t) = k x(t) D x(t) , (4.78)

m

where

! !

(ka + kb )/m kb /m !2 + !2 !2b

D= = a 2 b (4.79)

kb /m kb /m !b !2b

is the dynamical matrix with !a = ka /m and !b = kb /m. Equation (4.78), x(t) = D x(t),

together with the definitions of D and x, is just a compact way of rewriting Eqs. (4.75) and

(4.76). It is very useful, however, because it permits us to use the mathematical machinery

of linear algebra to solve these equations.

We pause here to take note of an important mathematical property of the stiness matrix

k : it is a symmetric matrix, meaning that ki j = k ji , where ki j are the components of k . We

will not go into the reason why just yet, but it is a quite general property of the stiness

matrix (which you can use as a check on your algebra: if the stiness matrix you have

found is not symmetric, then you have made a mistake!). The dynamical matrix D is also

symmetric, trivially so here as it is the same as k aside from the scalar multiplicative factor

of 1/m. We will make use of this symmetry property shortly.

To solve Eq. (4.78), we follow the same procedure we used for solving problems with

single oscillators; we look for solutions of the form

i!t

x = ae , (4.80)

where

!

a1

a= . (4.81)

a2

Substituting Eq. (4.80) into Eq. (4.78) yields

!2 a = D a , (4.82)

which, upon rearranging can be written in the form of an eigenvalue equation

. (4.83)

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92 Normal Modes

where = !2 is called the eigenvalue and a is known as the eigenvector. Our task is to

determine the eigenvalues and associated eigenvectors a for the matrix D , which contain

all the essential physics of the problem. The eigenvectors will give us the normal frequen-

cies and the eigenvectors will give us the normal coordinates. In this case, there are two

eigenvalues and two eigenvectors, which follows from the fact that there are two degrees

of freedom, or equivalently that D is a 2 2 matrix.

Because D is a symmetric matrix, all its eigenvalues are guaranteed to be real. More-

over, while we do not prove

p it here, D is positive definite, meaning that its eigenvalues

are positive. Since ! = , the normal frequencies ! are real. Physically, this is what we

expect, as complex frequencies would imply damping, which isnt present in this problem.

To solve the eigenvalue problem, we rewrite Eq. (4.83) as Da = Ia, where I is the identity

matrix defined by Ii j = i j , and then subtract the right side from both sides of the equation:

(D I )a = 0 . (4.84)

Equation (4.84) has a solution only if the determinant of the matrix D I is zero, that is, if

det D I = 0 (see Appendix A). Writing out Eq. (4.84) explicitly for this problem gives

!

2 !2a + !2b !2b

det(D ! I ) = det =0, (4.85)

!2b !2b

which gives

(!2a + !2b )(!2b ) !4b = 0 . (4.86)

Expanding Eq. (4.86) gives:

2

(!2a + 2!2b ) + !2a !2b = 0 . (4.87)

The determinant of D is a second order polynomial called the characteristic polynomial

of the matrix D . We encountered this earlier in 4.1.1 (see Eq. (4.10)). In general, the

characteristic polynomial of an n n matrix is an nth -order polynomial. Here it is a second

order polynomial, which we can solve using the quadratic formula.

Instead of solving Eq. (4.87) for the general p case for arbitrary spring constants, we con-

sider the special case where kb /ka = !2b /!2a = 3/2, which simplifies the algebra, but still

serves to illustrate the method. With this choice, the solution to Eq. (4.87) is

0p 1 (

BB 3 + 1 CC 0.366 !2a

= !2 = BB@ 1CCA !2a . (4.88)

2 2.366 !2a

Thus, the two normal frequencies for this system are

s

p

3 1 p

!1 = !a 0.366 !a 0.605 !a (4.89)

2

s

p

3+3 p

!2 = !a 2.366 !a 1.538 !a . (4.90)

2

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93 Matrix formulation of normal modes

Associated with the normal frequencies !1 and !2 are their respective eigenvectors a1 and

a2 , which we can find using Eq. (4.84). To find

p the eigenvector ap1 associated with !1 , we

substitute a1 into Eq. (4.84) using 1 = !21 = 32 1 !2a and !2b = 23 !2a . This gives

! !

!2 + !2b !2b a1,1

D !21I a1 = a =0 (4.91)

!2b !2b a2,1

0 p 1 !

BB 3 !2 3 2C

2 !a C

CC a1,1

= BBB@ 2p3 a 2 1 2 A

C =0, (4.92)

a2,1

2 !a 2 !a

where the second subscript for a1,1 and a2,1 is a 1 to indicate that these are components of

the eigenvector for normal mode 1 (which has a normal frequency of !1 . Multiplying out

the top and bottom rows yields

p

3 2 3 2

!a a1,1 ! a2,1 = 0 (4.93)

2 2 a

p

3 2 1

!a a1,1 + !2a a2,1 = 0 . (4.94)

2 2

Solving either of these equations for a2,1 yields the same result:

p

a2,1 = 3 a1,1 . (4.95)

p

3+3

Repeating the same procedure but substituting a2 into Eq. (4.84) using 2 = !22 = 2 !2a

yields two similar equations. Keeping only the equation from the top row gives

p

1 2 3 2

! a1,2 ! a2,2 = 0 . (4.96)

2 a 2 a

which yields

1

a2,2 = p a1,2 . (4.97)

3

If a1 and a2 , are eigenvectors, then any constant times a1 and a2 is also an eigenvector

(which, fundamentally, is why the above procedure determines only the ratios a2,1 /a1,1 and

a2,2 /a1,2 ). We will denote the normalized (length 1) eigenvectors as e1 and e2 , which are

given by

0 1 1 0 p

3C

1

BBB 2 CCC BBB

2 C

CCC

e1 = BBB@ p CCCA , e2 = BBB@ C . (4.98)

3 1 A

2 2

What is the physical interpretation of these normal coordinates we have found? To answer

this question, we return to Eq. (4.80), x = a e i!t , and write it out explicitly for the first of

t

94 Normal Modes

! !

x1,1 a1,1 i!1 t i!1 t

= e = e1 Ae (4.99a)

x2,1 a2,1

0 1 0 1 1

BBe1,1 CC BBB 2 CCC

= BBB@ CCCA Ae i!1 t

= BBB@ p CCCA Ae i!1 t

, (4.99b)

e2,1 3

2

where the first subscript on xi, and ai, refers to the displacement and amplitude of mass 1

or mass 2, and the second subscript (after the comma) designates the normal mode (here,

normal mode 1). These equations tell us that in normal p mode 1, mass 1 and mass 2 al-

ways move with relative amplitudes of a2,1 /a1,1 = 3. The absolute overall scale of the

amplitude is determined by the initial conditions and set by the (complex) number A in

the above equations. The fact that both components of the eigenvector e1 , namely e1,1 and

e2,1 , are positive means that in normal mode 1, both particles are always displaced from

equilibrium in the same direction. Moreover, in mode 1, both particles oscillate with the

same frequency, !1 0.605 !a .

We can write similar equations for normal mode 2:

! !

x1,2 a1,2 i!2 t

= e = e1 Ae i!2 t (4.100a)

x2,2 a2,2

0 1 0 p3 1

BBBe1,2 CCC BBB CCC

= BB@ CCA Ae i! 2 t

= BBB@ 2 CCCA Ae i!2 t , (4.100b)

e2,2 1

2

p

In mode two a2,2 /a1,2 = 1/ 3, meaning that the two masses are always displaced in

opposite directions, albeit with dierent amplitudes. In mode 2, both particles oscillate

with the same frequency, !2 1.538 !a .

So to summarize, you should think of the normal coordinates as giving the relative am-

plitudes of all the particles in a given normal mode. All particles in that mode oscillate with

the same frequency, the normal frequency, also known as the eigenfrequency.

Figure 4.6 shows the eigenvectors e1 and e2 graphically, plotted in a basis (coordinate

system) defined by the particle coordinates (x1 , x2 ). Note that e1 and e2 are perpendicular

to each other and can be written as

0 1 0 1

BBBcos CCC BBB sin CCC

e1 = B@ B C

CA , e2 = B@ B CCA , (4.101)

sin cos

where = 60 for this problem. Other choices of the ratio kb /ka give similar results, with e1

and e2 perpendicular to each other and rotated with respect to e1 and e2 , but with dierent

angles . Note that we have chosen the signs for the components e1 and e2 to give a right

handed coordinate system. In particular, this meant choosing the upper component of the

e vector to be negative. Making such a choice is not strictly necessary, but it does make it

somewhat simpler to transform results between the dierent coordinate systems.

t

95 Matrix formulation of normal modes

t

Fig. 4.6 Normal coordinates for 2-mass 2-spring problem. Unit vectors e1 and e2 are shown

along the x1 and x2 directions (gray) and for the normal coordinates e1 and e2 (black).

The vectors e1 and e1 define a new coordinate system that is rotated by the angle with

respect to the original coordinate system defined by the unit vectors e1 and e2 along the

x1 and x2 axes, respectively. The solutions to the problem can be expressed in either set

of coordinates: in the e1 -e2 coordinate system, solutions are expressed as functions of the

original coordinates x1 and x2 ; in the e1 -e2 coordinate system, solutions are expressed as

functions of the rotated coordinates x1 and x2 . In the context of linear algebra, each of these

coordinate systems is referred to as a basis. Any vector x describing the positions of the

masses can be expressed either basis:

x = x1 e1 + x2 e2 = x1 e1 + x2 e2 . (4.102)

For example, we can express the components x1 and x2 in either basis using dot products:

x2 = x e2 = (x1 e1 + x2 e2 ) e2 = (e1 e2 ) x1 + (e2 e2 ) x2 . (4.104)

x = U x , (4.105)

where

! !

x1 x1

x = , x= , (4.106)

x2 x2

and

! !

e1 e1 e2 e1 cos sin

U= = (e1 , e2 ) = . (4.107)

e1 e2 e2 e2 sin cos

Note that the columns of U are the eigenvectors e1 and e2 we found in Eqs. (4.98) and

(4.101). The matrix U is a rotation matrix that transforms (rotates) coordinates x in the

original (e1 , e2 ) basis to coordinates x in the eigenvector (e1 , e2 ) basis.

The equation of motion in matrix form is given by

x = D x . (4.108)

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96 Normal Modes

We can rewrite Eq. (4.108) in the tilde basis. Noting from Eq. (4.105) that x = U 1 x, Eq.

(4.108) can be written as

U 1 x = DU 1 x . (4.109)

x = UDU 1 x , (4.110)

or

x = D

Dx , (4.111)

where

1

D UDU

D = UDU T , (4.112)

where we have noted the property of rotation matrices that their inverse is their transform.

The two eigenvectors, with eigenvalues 1 = !21 and 2 = !22 , point along the e1 and e2

axes, respectively, so in the tilde basis,

! !

1 0

e1 = , e2 = . (4.113)

0 1

D a = a, D

Since D D must be a diagonal matrix with the eigenvalues as the diagonal elements

! !

1 0 !21 0

D UDUT = = . (4.114)

0 2 0 !22

obtain two uncoupled equations of motion

x = D

Dx

x e = D

D x e = x e ) x = !2 x , (4.115)

with = 1, 2. Thus, our procedure has reduced the original two coupled equations of

motion to two uncoupled equations of motion.

To find solutions for a specific set of initial conditions, we proceed the same way we pro-

ceeded in Example 4.1, remembering that we are now using the notation x in place of q to

refer to the normal coordinates. Thus, as before, since x1 and x2 obey equations of motion

for uncoupled simple harmonic oscillators, the solutions can be written in any one of the

usual forms. Here we use sines and cosines

x2 (t) = A2 cos !2 t + B2 sin !2 t , (4.117)

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97 Matrix formulation of normal modes

and determine the four amplitudes from the initial conditions, which are generally ex-

pressed in terms of the original variables. For example, suppose x1 (0) = 4 and x2 (0) = 2

and x1 (0) = x2 (0) = 0, which means that

! !

4 0

x= , x = . (4.118)

2 0

We can use Eq. (4.105) to find the initial conditions for the normal coordinates

! ! !

x1 (0) cos sin x1 (0)

x(0) = = Ux = (4.119)

x2 (0) sin cos x2 (0)

0 p 10 1

3C

0 p 1 0 1

BBB 1 CCC BBB4CCC BBB 2 + 3 CCC BBB 3.732 CCC

B

= BB@B p

2 2 B

CCC BB@ CCA = BB@ C

p CCA ' BB@ CCA , (4.120)

3 1 A 2 1 2 3 2.464

2 2

! ! !

x1 (0) cos sin x1 (0)

x (0) = = U x = (4.121)

x2 (0) sin cos x2 (0)

0 p 10 1 0 1

BBB 1 3C

CC BBB0CCC BBB0CCC

2 C

= BBBB@ p

2 CCC BB@ CCA = BB@ CCA , (4.122)

3 1 A 0 0

2 2

Applying these initial conditions to Eqs. (4.116) and (4.117) gives A1 ' 3.732, B1 = 0,

A2 ' 2.464, and B2 = 0. To find solutions for x1 (t) and x2 (t), we use the inverse of Eq.

(4.105)

0 1 0 p 1

3C 0 1

BBB x1 (t)CCC BBBB 12 CC BBB x1 (t)CCC

2 C

BB@ CCA = BBB p CCC B@B CAC , (4.124)

x2 (t) @ 3 1 A x (t)

2

2 2

1

where we have used the property of rotation matrices that U = U T . With these results,

we obtain

p

3

x1 (t) = 12 x1 (t) 2 x2 (t) (4.125)

p

1 3

' 2 3.732 cos !1 t 2 ( 2.464 cos !2 t) (4.126)

' 1.866 cos !1 t + 2.1344 cos !2 t . (4.127)

and

p

3 1

x2 (t) = 2 x1 (t) + 2 x2 (t) (4.128)

p

3

' 2 3.732 cos !1 t + 12 ( 2.464 cos !2 t) (4.129)

' 3.232 cos !1 t 1.232 cos !2 t . (4.130)

t

98 Normal Modes

We would like to find the normal modes for a system consisting of three equal masses

attached to a massless string stretched between two fixed walls, as depicted in Fig. 4.7.

The string is held under tension T with the masses equally spaced at distance l from each

other and from the walls. Suppose one or more of the masses is pulled up and then released

setting the system into oscillation in the vertical plane. The vertical distance of each mass

from its equilibrium position is yi , where i = 1, 2, 3. We take the amplitude of the oscilla-

tions to be small compared to the distance between masses so that yi l at all times.

The first step is to write down the equations of motion for each mass. Lets begin with

the leftmost mass whose vertical position is y1 . The vertical force on the mass will be the

dierence in the vertical component of the tension of the string on either side of the mass

In the limit of small amplitude oscillations where yi l, the sines can be approximated as

sin 0 ' y1 /l and sin 1 ' (y2 y1 )/l. Thus, the force on the leftmost mass is

y1 y2 y1 T

F1y = T +T = ( 2y1 + y2 ) (4.132)

l l l

Using similar analysis, the vertical forces on the middle and rightmost masses are

y2 y1 y3 y2 T

F2y = T +T = (y1 2y2 + y3 ) (4.133)

l l l

y2 y3 y3 T

F3y =T T = (y2 2y3 ) (4.134)

l l l

You should check that the signs of the vertical forces are in the correct directions for the

configurations of the masses used in Fig. 4.7. We can now write down the equations of

motion of each of the three masses:

T

my1 = ( 2y1 + y2 ) (4.135)

l

T

my2 = (y1 2y2 + y3 ) (4.136)

l

T

my3 = (y2 2y3 ) . (4.137)

l

t

Fig. 4.7 Three masses on a string under tension. The vertical scale is greatly magnified

compared to the horizontal.

t

99 Matrix formulation of normal modes

0 1

BBB 2 1 0 CCC

B C

D = !20 BBBB 1 2 1CCCCC , (4.138)

B@ A

0 1 2

p

where !0 = T/ml. The characteristic frequency !0 that emerges increase when the ten-

sion in the string is increased or when the mass or distanced between the masses is in-

creased, which makes sense physically.

The normal frequencies are found by finding the solutions to the characteristic equation,

which is given by det(D I ) = 0 where = !2 . Defining s = /!20 , this becomes

0 1

BBB2 s 1 0 CCC

B C

det(D I ) = !20 det BBBB 1 2 s 1 CCCCC = 0 (4.139)

B@ A

0 1 2 s

h i

= !20 (2 s)3 2(2 s) = 0 , (4.140)

p

which has solutions of s = 2 and s = 2 2. Therefore, the normal frequencies, from the

lowest to the highest frequency, are given by

q

p

!1 = 2 2 !0 0.765 !0 (4.141)

p

!2 = 2 !0 1.414 !0 (4.142)

q

p

!3 = 2 + 2 !0 1.848 !0 , (4.143)

p

where !0 = T/ml.

The eigenvectors are found by solving (D !2I )a = 0 for each normal frequency !

with = 1, 2, 3. Starting with !1 , we have (D !21I )a1 = 0, which written out gives

0 p 10 1

BBB2 (2 2) 1 0 CCC BBBa1,1 CCC

BBB p CCC BBB CCC

BBB 1 2 (2 2) 1

p CCC BBBa2,1 CCC = 0 (4.144)

@ A@ A

0 1 2 (2 2) a3,1

Multiplying out the three rows, we obtain

p

2 a1,1 a2,1 = 0 (4.145)

p

a1,1 + 2 a2,1 a3,1 = 0 (4.146)

p

a2,1 + 2 a3,1 = 0 . (4.147)

p p

From the first and third equations we have a2,1 = 2 a1,1 and a2,1 = 2 a3,1 , which means

that a3,1 = a1,1 . The second equation is automatically satisfied by these relations. Thus, the

first normalized eigenvector is

0 1

B 1 C

1 BBBBB p CCCCC

e1 = BB 2CC . (4.148)

2 B@ CA

1

t

100 Normal Modes

0 10 1

BBB2 2 1 0 CCC BBBa1,2 CCC

BBB CB C

BBB 1 2 2 1 CCCCC BBBBBa2,2 CCCCC = 0 (4.149)

@ A@ A

0 1 2 2 a3,2

Again, multiplying out the three rows, we obtain

a2,2 = 0 (4.150)

a1,2 a3,2 = 0 (4.151)

a2,2 = 0 . (4.152)

The first and third equations give a2,2 = 0. The second equation gives a1,2 = a3,2 . Thus,

the second normalized eigenvector is

0 1

B1C

1 BBBBB CCCCC

e2 = p BB 0 CC . (4.153)

2 B@ 1CA

0 p 10 1

BBB2 (2 + 2) 1 0 CCC BBBa1,3 CCC

BBB p CCC BBB CCC

BBB 1 2 (2 + 2) 1 a =0 (4.154)

@ p CCCA BBB@ 2,3 CCCA

0 1 2 (2 + 2) a3,3

Multiplying out the three rows, we obtain

p

2 a1,1 a2,1 = 0 (4.155)

p

a1,1 2 a2,1 a3,1 = 0 (4.156)

p

a2,1 2 a3,1 = 0 . (4.157)

p p

From the first and third equations we have a2,1 = 2 a1,1 and a2,1 = 2 a3,1 , which

means that a3,1 = a1,1 . The second equation is automatically satisfied by these relations.

Thus, the third normalized eigenvector is

0 1

BBB 1 CCC

1B p C

e3 = BBBB 2CCCCC . (4.158)

2 B@ A

1

Figure 4.8 shows the normal modes for three masses on a string under tension. The three

components of each of the three eigenvectors give the relative amplitudes of the masses in

each normal mode.

Example 4.2 Consider the case of three masses on a string where m = 25 g, l = 10 cm,

and T = 50 N. Suppose that the middle mass is given an initial velocity of 30 = 0.50 m/s,

while the other two masses are initially stationary. All three masses start out in their equi-

librium positions. What is the subsequent motion of the masses? What frequencies are

produced?

t

101 Matrix formulation of normal modes

t

Fig. 4.8 Normal modes of 3 masses on a string. The relative amplitudes in each normal mode

are proportional to the components of the corresponding eigenvector. The

eigenvectors and normal frequencies are shown to the left and right, respectively, of

each drawing.

Solution

Lets start by calculating the characteristic vibration frequency !0 :

r s

T 50 N 1

!0 = = = 447 s (4.159)

ml (0.025 kg)(0.010 m)

!0

f0 = = 71.2 Hz . (4.160)

2

The initial conditions are

Note that with these initial conditions, masses 1 and 3 must have identical trajectories,

because they are symmetrically placed about mass 2, which is the only mass that is initially

perturbed from equilibrium. This means that normal mode 2 will not be excited in this case,

because in normal mode 2, masses 1 and 3 are displaced in opposite directions. Another

way of looking at this is to note that the initial conditions are symmetric with respect to the

central particle, as are normal modes 1 and 3. By contrast, normal mode 2 is antisymmetric

with respect to mass 2, so it will not be excited in this case.

Generalizing Eqs. (4.103) and (4.104) to a 3-particle system, the normal coordinates are

1h p i

y1 (t) = y e1 = y1 (t) + 2 y2 (t) + y3 (t) (4.162)

2

1

y2 (t) = y e2 = y1 (t) y3 (t) (4.163)

2

1h p i

y3 (t) = y e3 = y1 (t) 2 y2 (t) + y3 (t) . (4.164)

2

t

102 Normal Modes

y1 = y2 = y3 = 0 (4.165)

1h p i 1

y 1 = y1 (0) + 2 y2 (0) + y3 (0) = p 30 (4.166)

2 2

1

y 2 = p y1 (0) y3 (0) = 0 (4.167)

2

1 h p i 1

y 3 = y1 (0) 2 y2 (0) + y3 (0) = p 30 . (4.168)

2 2

We will write the solutions to each normal mode as a sum of a sine and a cosine function.

The general solutions for the dierent normal modes are

y2 (t) = A2 cos !2 t + B2 sin !2 t (4.170)

y3 (t) = A3 cos !3 t + B3 sin !3 t . (4.171)

verifying our earlier claim that normal mode 2 would not appear in the solution for this set

of initial conditions. Applying the remaining initial conditions

1

y 1 = !1 B1 = p 30 (4.172)

2

1 30 1 30

) B1 = p = q p !0 = 1.03 mm (4.173)

2 !1 2 2 2

1

y 3 = !3 B3 = p 30 (4.174)

2

1 30 1 30

) B3 = p = q p = 2.92 mm . (4.175)

2 !3 2 2+ 2 0

!

Note that B1 l and B3 l, justifying our approximation that the deflections of the string

are small compared to the spacing between particles. Because y2 (t) = 0, y1 (t) = y3 (t).

Adding and subtracting Eqs. (4.162) and (4.164) thus gives

1 1

y1 (t) = y3 (t) = y1 (t) + y3 (t) = [B1 sin !1 t + B3 sin !3 t] (4.176)

2 2

1 1

y2 (t) = p y1 (t) y3 (t) = p [B1 sin !1 t B3 sin !3 t] . (4.177)

2 2

The analytical method we have developed thus far is suited only for coupled systems of

equal masses. When the masses are not equal, forming the dynamical matrix becomes a bit

more involved.

Lets consider the problem of two pendulums coupled by a spring with a force constant k

t

103 Matrix formulation of normal modes

that we introduced at the beginning of this chapter, but this time lets make the two masses,

denoted M and m, dierent. Figure 4.1 still serves to define the problem and aside from the

dierent masses, the equations of motion are the same as they were before

mg

m x1 = + k x1 + kx2 (4.178)

l

Mg

M x2 = kx1 + k x2 , (4.179)

l

We can write these equations in matrix form

m x = k x , (4.180)

where

! ! !

m 0 mg/l + k k x1

m= , k= , x= . (4.181)

0 M k Mg/l + k x2

We wish to recast Eq. (4.180) in the standard form x = D x. The most straightforward

approach is simply to left multiply both sides of Eq. (4.180) by the inverse of the mass

matrix m 1 , which gives

x = D 0 x (4.182)

where

! !

m 1 0 mg/l + k k

D0 = m 1k = (4.183)

0 M 1 k Mg/l + k

! !

g/l + k/m k/m !2 + !2 !2s

= = p 2s , (4.184)

k/M g/l + k/M ! s !2p + !2s

p p

where ! p = g/l, ! s = k/m, and = m/M. The eigenvalues of D 0 are given by

solving det(D 0 I ) = 0:

!

!2p + !2s !2s

det = (!2p + !2s )(!2p + !2s ) !4s = 0 (4.185)

!2s !2p + !2s

=( !2p )( !2p (1 + )!2s ); . (4.186)

Solving for = !2 gives

8

>

> 2

<! p

2

=! =>

> , (4.187)

:!2 + (1 + )!2

p s

which reduces to the result we obtained in 4.1.1 for equal masses when = 1.

We could go on to find the eigenvectors, which proceeds along now familiar lines. How-

ever, the dynamical matrix D 0 is not symmetric, and this has certain disadvantages. First,

when the dynamical matrix is not symmetric, the eigenvectors do not in general form an

orthogonal system. As we shall see, the can be a disadvantage in performing certain calcu-

lations. Second, for systems with a large number of degrees of freedom, the eigenvalues and

eigenvectors must be found numerically in most cases. Because a non-symmetric matrix

has nearly twice as many distinct entries as its symmetric counterpart for which Di j = D ji ,

t

104 Normal Modes

numerical computation of eigenvalues and eigenvectors require more memory and take

considerably more time for a non-symmetric matrix than for its symmetric counterpart.

For the diagonal mass matrices m we have encountered here, it is a simple matter to

rewrite the eigenvalue equation so that the dynamical matrix is diagonal. We start by defin-

ing the square root m 1/2 and inverse square root m 1/2 of the mass matrix

! !

1/2 m1/2 0 1/2 m 1/2 0

m = , m = . (4.188)

0 M 1/2 0 M 1/2

1/2

Starting from m x = k x, Eq. (4.180), we left multiply both sides by m

m 1/2m x = m 1/2k m 1/2m 1/2 x , (4.189)

1/2 m 1/2

and insert the identity matrix I = m m as shown on the right hand side. Regrouping

terms, this becomes

m 1/2 x = m 1/2

km 1/2

m 1/2 x , (4.190)

or

y = D y , (4.191)

where

p ! p !

1/2 1/2 g/l + k/m k/ mM !2 + !2 !2s

D =m km = p = pp 2s , (4.192)

k/ mM g/l + k/M !s !2p + !2s

is now a symmetric dynamical matrix and

! ! p !

m1/2 0 x1 m x1

y = m 1/2 x = = p . (4.193)

0 M 1/2 x2 M x2

The only price we have to pay for this transformation is the introduction of the vector y. A

quick examination of D 0 and D reveals that they have the same characteristic polynomial,

and thus the same set of eigenvalues. However, they do not have the same eigenvectors.

To find the eigenvectors, we need to solve (D !2I )a = 0 for each frequency ! . While

we can perform the calculation for arbitrary ! p and ! s , the algebra is messy. To simplify

the algebra, we choose ! p = 6, ! s = 2, and = m/M = 9/4. For this case, Eq. (4.192)

gives

!

40 6

D= (4.194)

6 45

while Eq. (4.187) gives = 36 and 49 or !1 = 6 and !2 = 7. The eigenvectors are given

by solving (D !2I )a = 0 for !1 = 6 and !2 = 7:

! ! ! !

40 36 6 a1,1 40 49 6 a1,2

=0, =0, (4.195)

6 45 36 a2,1 6 45 49 a2,2

which gives normalized eigenvectors of

! ! ! !

1 3 cos 1 2 sin

e1 = p = , e2 = p = , (4.196)

13 2 sin 13 3 cos

t

105 Matrix formulation of normal modes

where as before the the normalized two-component eigenvectors can be written in terms of

sines and cosines of the rotation angle ; in this case ' 33.7 .

We can now use these results to solve an initial value problem. Lets take initial condi-

tions of x1 (0) = x0 , x2 (0) = 0, and x1 (0) = x2 (0) = 0. We need to express these in terms of

y(0) and y(0), which we do using Eq. (4.193)

! p ! p !

y1 (0) m x1 (0) m x0

y(0) = = p = (4.197)

y2 (0) M x2 (0) 0

! p ! !

y1 (0) m x1 (0) 0

y(0) = = p = . (4.198)

y2 (0) M x2 (0) 0

The general solutions, written in terms of the normal coordinates y1 (t) and y2 (t), are given

by the usual sine and cosine functions

1

y1 (t) = p [3y1 (t) + 2y2 (t)] = A1 cos !1 t + B1 sin !1 t (4.199)

13

1

y2 (t) = p [ 2y1 (t) + 3y2 (t)] = A2 cos !2 t + B2 sin !2 t . (4.200)

13

The initial conditions on y1 and y2 are

r

1 m

y1 (0) = p 3y1 (0) + 2y2 (0) = 3x0 (4.201)

13 13

r

1 m

y2 (0) = p 2y1 (0) + 3y2 (0) = 2x0 (4.202)

13 13

y 1 (0) = y 2 (0) = 0 . (4.203)

The initial conditions on the velocities y 1 (0) and y 2 (0) give B1 = B2 = 0. Thus, the solu-

tions are given by

r r

m m

y1 (t) = 3x0 cos !1 t , y2 (t) = 2x0 cos !2 t . (4.204)

13 13

p

Solving Eqs. (4.199)

p and (4.200) for y1 and y2 , we obtain y1 (t) = (1/ 13)[3y1 (t) 2y2 (t)]

and y2 (t) = (1/ 13)[2y1 (t) + 3y2 (t)]. Multiplying Eq. (4.193) by m 1/2 gives

0 1 1 0 1 1

BBB pm y1 CCC BBB p [3y1 (t) 2y2 (t)] CCC

13m

x = m 1/2 y = BBB@ 1 CCCA = BBB@ 1 CCCA . (4.205)

p y2 p [2y1 (t) + 3y2 (t)]

M 13M

Substituting the expressions for y1 (t) and y2 (t) from Eq. (4.204) into Eq. (4.205) gives

!

9 4

x1 (t) = x0 cos !1 t + cos !2 t (4.206)

13 13

p

6

x2 (t) = x0 (cos !1 t cos !2 t) . (4.207)

13

A quick check on these solutions reveals that they satisfy the initial conditions that x1 (0) =

x0 , x2 (0) = 0, and x1 (0) = x2 (0) = 0. They also satisfy the equations of motion, as must

any solution that is a sum sine and cosine functions of !1 t and !2 t.

t

106 Normal Modes

x1 (t)

1.0

0.5

0.0 t

2 4 6 8 10 12 14

0.5

1.0

x2 (t)

1.5

1.0

0.5

0.0 t

2 4 6 8 10 12 14

0.5

1.0

t

1.5

masses. The dashed lines show that time-dependent amplitudes of the oscillations.

When !1 and !2 are nearly equal, we expect beats to appear. In this case, it is useful to

rewrite the solutions in terms of the average and dierence frequencies, ! 12 (!1 + !2 )

and ! 12 (!2 !1 ), respectively. The normal mode frequencies are then written as

!1 = ! ! and !2 = ! + !. Using cosine addition identities, the solutions can be

rewritten as

!

5

x1 (t) = x0 cos !t cos !t

+ sin !t sin !t

(4.208)

13

p

12

x2 (t) = x0 sin !t sin !t

(4.209)

13

The trajectories for the two pendulums are shown in Fig. 4.9. As in the case treated

previously for equal-mass coupled pendulums, a beating pattern is observed. In contrast

to the previous case, however, the amplitude of the oscillations of the heavier mass goes

through a minimum but never goes to zero. More generally, one can show3 that

+ (cos2

x1 = x0 cos !t cos !t sin2 ) sin !t sin !t

, (4.210)

which means that the case where x1 (t) goes to zero when x2 (t) is a maximum is a special

case that only occurs when cos2 sin2 = 0 as is does, for example, when = 45 . This

is precisely what occurs for the case when m = M, which we studied in 4.1.1.

3

p

To obtain

p the more general result, redo the calculation of starting with Eq. (4.199) and let (2/ 13) ! sin

and (3/ 13) ! cos .

t

107 Matrix formulation of normal modes

t

Fig. 4.10 Coordinates for a coupled pendulum with unequal masses. (a) Normal coordinates e01

in (x1 , x2 ) coordinates: ' 69 and 0 = 45 . (b) Normal coordinates in stretched

(y1 , y2 ) coordinates: = 90 and ' 33.7 .

The transformation to a coordinate system in which D is symmetric may seem like a lot

more trouble than it is worth. In some cases that may be true, but there are practical ben-

efits that may not be readily apparent. The most important one is that the eigenvectors are

orthogonal when D is symmetric, and many other useful properties follow from this one.

Lets take a closer look at the problem of the pendulum with unequal masses that we

just solved. Once again we consider the particular case in which ! p = 6, ! s = 2, and

= m/M = 9/4, and calculate the eigenvectors of D0 in the original basis. In that basis, the

dynamical matrix is given by

!

40 4

D0 = , (4.211)

9 45

which is not symmetric. The eigenvalues of D 0 are = 36 and 49 or !1 = 6 and !2 = 7,

the same as for D , as they must be. But the eigenvectors are dierent. With a little algebra,

the normalized eigenvectors of D 0 are found to be

! !

0 1 1 0 1 4

e1 = p , e2 = p . (4.212)

2 1 97 9

The cosine of the angle 0 between e01 and e02 is given by their dot product, from which we

find that ' 69 : the eigenvectors are not orthogonal. Figure 4.10(a) shows the geometry

of the eigenvectors e01 and e01 .

Figure 4.10(b) illustrates how the transformation to the y basis, given by Eq. (4.193),

stretches the (x1 , x2 ) coordinate system so that the basis vectors become orthogonal, with

= 90 in the new basis. After stretching the basis to obtain orthogonal normal coordi-

nates, the coordinates are renormalized to have unit length once again.

Thus, while the ostensible purpose of the transformation to the y basis was to obtain a

symmetric dynamical matrix, from the standpoint of geometry, the purpose of the trans-

formation is to find a coordinate system in which the eigenvectors are orthogonal. It is this

orthogonality that makes it possible to obtain the normal coordinates as a pure rotation and

to express the eigenvectors in terms of sines and cosines of the rotation angle. This, in turn,

allows us to obtain the result given by Eq. (4.210), which helps us understand how beats

t

108 Normal Modes

can appear either fully modulated, as seen in Fig. 4.4, or only partially modulated, as seen

in Fig. 4.9.

Formulating normal modes in the language of linear algebra and matrices brings a little

more order to the study of the subject. It also allows us to draw on a vast trove of math-

ematical theorems from linear algebra, which simplifies the analysis of such problems, as

well as giving physical insight. Thus far, we have only studied systems with a few degrees

of freedom. In Chapter 5 we take up the problem of normal modes again and apply matrix

methods to systems with a large number of degrees of freedom. For systems with many de-

grees of freedom, transforming the coordinates to make the dynamical matrix symmetric is

particularly important, as it is generally much more efficient computationally for systems

with many degrees of freedom.

Carbon dioxide is a linear molecule consisting of a single carbon atom sandwiched be-

tween two oxygen atoms. If we consider only small oscillations about the equilibrium

configurations of the three atoms, we can treat the problem of molecular vibrations within

the framework of normal modes that we have developed in this chapter.

To begin, we consider only longitudinal motion, that is, motion along the axis formed by

the three atoms. Later we will consider tranverse motion, which is motion perpendicular to

the molecular axis. The longitudinal deviations of the three particles from their equilibrium

positions are denoted x1 , x2 , and x3 , as shown in Fig. 4.11(a). The mass of the two oxygen

atoms is m and the mass of the carbon atom is M. The equations of motion for the three

particles are

F1 = k(x2 x1 ) = m x1 (4.213)

F2 = k(x3 x2 ) k(x2 x1 ) = M x2 (4.214)

F3 = k(x2 x3 ) = m x3 , (4.215)

where the spring constant is given by

d2 U

k= , (4.216)

d2 =0

U() is the interatomic potential between the oxygens and the carbon atom and = x2 x1

or = x3 x2 . Following the usual procedure and writing xi (t) = ai e i!t for i = 1, 2, 3, we

obtain Eq. (4.78), which we repeat here for easy reference,

m x = k x , (4.217)

with the mass and stiness matrices obtained from Eqs. (4.213)-(4.215) given by

t

109 Normal modes of carbon dioxide

t

Fig. 4.11 Geometry for vibrations of triatomic CO2 molecule: (a) longitudinal, (b) transverse.

0 1 0 1

BBBm 0 0 CCC BBB k k 0 CCC

B C B C

m = BBBB 0 M 0 CCCCC , k = BBBB k 2k kCCCCC . (4.218)

B@ A B@ A

0 0 m 0 k k

The inverse of m is

0 1 1

BBBm 0 0 CCC

B C

m 1

= BBBBB 0 M 1 0 CCCCC , (4.219)

@ A

0 0 m 1

so that the dynamical matrix is given by

0 1

BBB k/m k/m 0 CCC

B C

D = m 1k = BBBB k/M 2k/M k/M CCCCC , (4.220)

B@ A

0 k/m k/m

which although not symmetric, will be used as is. Once again, the solution to the systems of

equations is given by solving the secular equation det(D I ) = 0, which for this problem

is

0 1

BBBBk/m k/m 0 CCC

C

det(D I ) = det BBBB k/M 2k/M k/M CCCCC = 0 , (4.221)

B@ A

0 k/m k/m

!" ! ! #

k 2k k 2k2

=0, (4.222)

m M m mM

which upon simplification yields

!" !#

k k 2m

1+ =0. (4.223)

m m M

Therefore, the normal frequencies are

r s !

k k 2m

! = 0 ! = ! = 1+ , (4.224)

m m M

which corresponds to three normal modes and six solutions ! , ! , and ! . At first

glance, the slowest mode with ! = 0 might seem to be a mistake, until we recall that the

equation of motion for a normal mode is always of the form

q + !2 q = 0 . (4.225)

t

110 Normal Modes

q = 0 , (4.226)

which has the general solution q (t) = q (0) + 3 t, or motion along a straight line an

constant velocity 3 . Thus, one of our normal modes is not a vibrational mode at all but

corresponds to constant velocity rigid body motion of the center of mass.

Next we proceed to obtain the eigenvectors for each normal mode. This will reveal the

nature of the each of the modes corresponding to the normal frequencies ! , ! , and ! .

Starting with the eigenvalue equation Eq. (4.83)

Da = a,

we set = !2 = 0 to examine the normal mode. For this case, Eq. (4.83) becomes

0 10 1 0 1

BBB k/m k/m 0 CCC BBBa1 CCC BBB0CCC

BBB CB C B C

BBB k/M 2k/M k/M CCCCC BBBBBa2 CCCCC = BBBBB0CCCCC . (4.227)

@ A@ A @ A

0 k/m k/m a3 0

Multiplying out the top row gives a1 a2 = 0, which means that a1 = a2 . Multiplying

out the bottom row gives a3 a2 = 0, which means that a3 = a2 . Thus, the normalized

eigenvector for the mode is aT = p13 (1, 1, 1), which is consistent with our earlier finding

that all three particles move together in a zero-frequency mode corresponding to rigid

translation.

To find the second normal mode, we set = !2 = k/m. Equation (4.83) becomes

0 10 1 0 1

BBB k/m k/m 0 CCC BBBa1 CCC Ba C

BBB CCC BBB CCC k BBBBB 1 CCCCC

BBB k/M 2k/M k/M CCC BBBa2 CCC = BBBa2 CCC . (4.228)

@ A@ A m @ A

0 k/m k/m a3 a3

Multiplying out the top row gives a1 a2 = a1 , which means that a2 = 0. Setting a2 = 0 and

multiplying out the middle row gives a3 = a1 . Thus, the normalized eigenvector for the

mode is aT = p12 (1, 0, 1). This means that the mode is one in which the two oxygen

atoms move exactly out of phase with each other and the central carbon atom remains

fixed, as does the center of mass of the molecule. Because the carbon atom does not move,

p

its mass M does not figure in the normal frequency ! = k/m for this problem.

To find the third normal mode, we set = !2 = (k/m)(1 + 2m/M). Equation (4.83)

becomes

0 10 1 0 1

BBB k/m k/m 0 CCC BBBa1 CCC " !# BBa1 CC

BBBB k/M 2k/M CB C k 2m BBBB CCCC

BB@ k/M CCCCC BBBBBa2 CCCCC = 1+ Ba2 C . (4.229)

A@ A m M BB@ CCA

0 k/m k/m a3 a3

Multiplying out the top row gives a2 = (2m/M)a1 . Similarly, multiplying out the bottom

row gives a2 = (2m/M)a3 , which means p that a1 = a3 . Thus the normalized eigenvector

T

for the mode is a = (1, 2m/M, 1)/ 2 + 4(m/M)2 ).

To summarize, the three normal modes are characterized by the following eigenvectors

t

111 Normal modes of carbon dioxide

0 1

B1C

1 BBBBB CCCCC

a = p BB1CC , !2 = 0 (4.230)

3 B@1CA

0 1

B 1C

1 BBBBB CCCCC k

a = p BB 0 CC , !2 = (4.231)

2 B@ 1 CA m

0 1

BBB 1 CCC !

1 BBB C k 2m

a = p BBB 2m/M CCCCC , !2 = 1+ (4.232)

2

2 + 4(m/M) ) @ A m M

1

Figure 4.12(a) depicts the displacements of the atoms for each normal mode. Note that the

relative amplitudes of the displacements correspond to the relative amplitudes of the three

components of the corresponding eigenvector. This is a general feature of the eigenvectors

and thus provides a simple connection between the components of components of

The reason we obtained three normal modes from our analysis is because we allowed

for three degrees of freedom. We could have eliminated the trivial zero-frequency ! mode

from the analysis at the outset by setting up the problem such that the center of mass of the

system is fixed. The center of mass (along the x axis) of the system is

1 X 1

xcm = P mi xi = [m(x1 + x3 ) + Mx2 ] . (4.233)

mi 2m + M

By choosing our coordinate system such that xcm = 0, we can eliminate one of the coordi-

nates, say x2 , by writing

m

x2 = (x1 + x3 ) . (4.234)

M

Substituting this back into Eqs. (4.213)-(4.215) yields two independent equations of mo-

tion, which can be solved in the usual fashion with the result that only the normal frequen-

cies corresponding to ! and ! appear as solutions. We leave this as an exercise at the

end of the chapter.

t

Fig. 4.12 Displacements for normal modes of CO2 : (a) longitudinal modes , , and , along the

x direction, (b) transverse mode along the y direction. The mode is not a

vibrational mode but corresponds to pure translation. Other pure translational and

rotational modes are not shown.

t

112 Normal Modes

Thus far, we have only allowed for motion along a line colinear with the equilibrium po-

sitions of the atoms, i.e. along the x axis. The modes thus obtained are called longitudinal

modes. There also exist transverse modes, that is, modes perpendicular to the x axis, which

we can study by allowing for motion in the y and z directions. However, as we have seen

from our analysis of the longitudinal modes, the number of modes we obtain can depend

on how we set up the problem. Before diving into another calculation, perhaps it is best

to think a bit. In particular, it would be useful to know how many transverse modes we

are looking for. On quite general grounds, we know that the number of modes is equal to

the number of degrees of freedom. Since each particle has three degrees of freedom and

there are three particles, there are a total of 3 3 = 9 degrees of freedom. Lets count the

dierent types. There are 3 translational degrees of freedom in which the whole molecule

moves in a straight line along the x, y, or z direction. There are 2 rotational degrees of free-

dom: the molecule can rotate about the y or z axis (but not about the x axis). In the previous

section, we found 2 longitudinal degrees of freedom. That leaves us with 9 3 2 2 = 2

transverse degrees of freedom. Based on the linear symmetry of the CO2 molecule, we

know that for any mode we find that involves motion only in the y direction, there must be

a corresponding completely independent mode oscillating at the same frequency in the z

direction. When two modes such as these have the same frequency, we say they are degen-

erate.4 These 2 degenerate transverse modes account for the last 2 modes of the system.

Since the two modes are degenerate, we have only one more mode to find.

Transverse y modes

We begin by determining the center of mass in the y direction and setting it equal to zero

in order to preclude the sort of trivial translational mode we encountered in the previous

section:

1 X 1

ycm = P mi yi = [m(y1 + y3 ) + My2 ] . (4.235)

mi 2m + M

m

y2 = (y1 + y3 ) . (4.236)

M

The symmetry of having two identical masses on either side of the central mass suggests

normal modes where either: (1) the two outer masses move in opposite directions along

the y axis and the central mass remains still so that ycm = 0, or (2) the two outer masses

move in the same direction along the y axis while the central mass moves in the opposite

direction so that ycm = 0. Indeed this is what we found for the longitudinal modes. Mode

(1) corresponds to the beginning of a rotational mode, which we have already counted and

excluded from this analysis. That leaves us with mode (2).

4 Degenerate modes are, in spite of their name, are no more reprobate than any other mode.

t

113 Normal modes of carbon dioxide

F1 = l = my1 (4.237)

F2 = 2l = M y2 (4.238)

F3 = l = my3 , (4.239)

where is the torsional spring constant for bending. Note that there is no obvious rela-

tionship between the torsional bending constant and the longitudinal spring constant k

introduced previously. The angle is related to the transverse displacements by

(y1 y2 ) + (y3 y2 )

= , (4.240)

l

in the limit of small displacements where yi l. For the mode we seek, y3 = y1 , yielding

2(y1 y2 )

' . (4.241)

l

Using Eq. (4.236) and noting once again that y3 = y1 , we obtain y2 = (2m/M)y1 , which

we can use to eliminate y1 from our expression for . This gives

2y2 M

' 1+ . (4.242)

l 2m

Substituting this expression for into Eq. (4.238) gives the equation of motion subject to

the constraints that ycm = 0 and y1 = y3 :

M

M y2 + 4 1 + y2 = 0 . (4.243)

2m

This is just the usual equation for a simple harmonic oscillator with an oscillation frequency

r

4 M

! = 1+ (4.244)

M 2m

The normalized eigenvector for the mode is

0 1

BBB 1 CCC

1 BBB C

a = p BBB 2m/M CCCCC (4.245)

2

2 + 4(m/M) ) @ A

1

(a) Write down the mass and stiness matrices m and k for the transverse modes

in terms of the coordinates (y1 , y2 , y3 ). You should find that both m and k are

symmetric.

(b) Calculate the matrix D = m 1k (which in this case is not symmetric) and verify

by direct substitution that (1, 2m/M, 1) is an (unnormalized) eigenvector with

a normal frequency given by Eq. (4.244).

(c) Verify that (1, 1, 1) and (1, 0, 1) are unnormalized eigenvectors with zero fre-

quency. Explain what each of these modes means physically. Does either mode

pose any problem for the yi 1 approximation we made?

t

114 Normal Modes

(d) Extra: Show that (2, 1, 0) is also an eigenvector with an eigenvalue of zero.

This mode cannot correspond to another independent normal mode because we

have already found the three independent transverse y modes that are allowed

for a system with three degrees of freedom. Show that the (2, 1, 0) is not an

independent mode.

In the previous section we found a single transverse normal mode in the y direction with

a nonzero frequency. Of course, that mode has an identical counterpart oscillating in the z

direction. The frequency of both modes is given by Eq. (4.244). The transverse oscillations

in the y direction are given by

y(t) = ey a (a y cos ! t + B y sin ! t) , (4.246)

where ey is a unit vector in the y direction, ! is given by Eq. (4.244), a is given by Eq.

(4.245), and the constants a y and B y are determined by the initial conditions. Similarly,

transverse oscillations in the z direction are given by

z(t) = ez a (a z cos ! t + B z sin ! t) . (4.247)

where the constants a z and B z are determined by the initial conditions.

Figure 4.13(a) shows the geometry for the two degenerate transverse modes. The direc-

tions y and z are arbitrary, of course, as long as we choose them to be perpendicular to the

molecular axis of the CO2 molecule. The oscillation frequency of these modes corresponds

to the oscillation frequency of infrared light with a wavelength of 15 m. If we illuminated

a container of CO2 with 15 m infrared radiation that was linearly polarized along a certain

direction, then it would excite the CO2 molecules to oscillate in that direction in a trans-

verse mode. In that case, it would be convenient to choose the y direction to correspond to

only one of the two degenerate modes. Otherwise, we would have to describe it as a linear

combination, the vector sum, of the two degenerate modes. There is nothing wrong with

doing this, other than making the problem a bit more complicated than necessary.

Suppose we were exciting the vibrational modes of CO2 with circularly polarized light.

Then we would expect that the atoms in the CO2 molecules would respond by vibrating in

one of the circularly polarized vibrational modes depicted in Fig. 4.13(b)-(c). Such modes

t

Fig. 4.13 Transverse modes of a carbon dioxide molecule: (a) y and z linear polarized modes,

(b) right circularly polarized, (c) left circularly polarized.

t

115 Damping and normal modes

can be readily created from linear combinations of Eqs. (4.246) and (4.247) by making the

two modes oscillate 90 out of phase with each other. For example, the left and right circu-

larly polarized modes shown in Fig. 4.13(b)-(c) can be realized by setting a y = B z = 1

and B y = a z = 0. In fact, any linear combination of the transverse modes is valid because

of the underlying linearity of the equations of motion.

If we are working with the circularly polarized modes, then it is a bit cumbersome to

describe them as phased superpositions of linearly polarized modes. It would be more

convenient if we could express them mathematically in a direct way. To facilitate this, we

first rewrite the linearly polarized modes using complex notation. Recalling that in so doing

we ultimately keep only the real part of our complex expressions, Eqs. (4.246) and (4.247)

become

i! t

z(t) = ez a a z e , (4.249)

where a y and a y are in general complex, and thus are able to carry the phase information

that is contained in the sine and cosine functions in Eqs. (4.246) and (4.247).

It is a simple matter to create circularly polarized modes from Eqs. (4.248) and (4.249)

above. Recalling that we keep only the real part, if we choose a z = ei/2 a y = ia y , the

we advance or delay the transverse z mode relative to the y mode, producing right and left

circularly polarized states, which we can write as

1

L(t) = p a a y ei! t ey + iez (4.250)

2

1

R(t) = p a a y ei! t ey iez , (4.251)

2

1/2

where the normalization factor of 2 assures that the column matrix has a length of unity.

The systems of coupled oscillators that we have examined thus far have been undamped.

Real systems exhibit damping, of course, so we proceed here to see what happens when

damping in included. The matrix formulation really comes in handy in this situation, as it

can teach us some general features about damped coupled oscillators that otherwise might

be difficult to figure out. To this end, we write down the equations of motion of a system

of damped coupled oscillators in matrix form:

m x = k x x . (4.252)

We proceed as before (see Eq. (4.80)) and substitute a trial solution of the form

i!t

x = ae , (4.253)

t

116 Normal Modes

which yields

!2m a = k a + i! a . (4.254)

1

Left multiplying both sides by m and rearranging terms so that the frequency ears only

on the right hand side gives

m 1k a = !2 a + i! m 1

a. (4.255)

Equation (4.255) is no longer an eigenvalue equation, which means that the normal mode

analysis we have developed cannot be used to find solutions. Nevertheless, Eq. (4.255) still

represents a system of N equations with N unknowns where N is the number of degrees

of freedom of the system of coupled oscillators. While it is possible to find the solutions

to these equations, the analysis can be tedious. So we will not attempt to find general

solutions to this equation. Instead we consider a special case for which Eq. (4.255) reduces

to an eigenvalue problem. Solutions for the more general case qualitatively resemble those

to the special case we are considering, and thus enable us to explore the basic physics of

the problem.

The special case we consider is the case when the damping matrix is proportional to the

mass matrix, that is, when

/m . (4.256)

1

In this case, m / I so that Eq. (4.255) can be rewritten as an eigenvalue equation,

Da = a , (4.257)

transforming Eq. (4.254), we would want to obtain a symmetric dynamical matrix, meaning

we would want to follow the procedure we developed in 4.2.5. Lets work through an

example.

Example 4.3 Consider the coupled pendulum we introduced at the beginning of this

chapter, but for the case where each pendulum experiences a damping force proportional

to its velocity. Assume the damping constants are the same for the two pendulums. Find

the normal modes, their frequencies, and the damping for each mode.

Solution

We start by writing down the equations of motion for the masses of the two pendulums,

including a damping term:

x1

F1 = mg + k(x2 x1 ) x1 = m x1 (4.258)

l

x2

F2 = mg k(x2 x1 ) x2 = m x2 , (4.259)

l

where we have included identical damping terms for each pendulum, and have used the

same notation and small angle approximations used previously to for the undamped case

t

117 Damping and normal modes

mg

m x1 = + k x1 + kx2 x1 (4.260)

l

mg

m x2 = kx1 + k x2 x2 , (4.261)

l

These two equations can be rewritten in matrix form following Eq. (4.252), with the mass

and stiness matrices given by Eq. (4.181) with M = m and the damping matrix given by

! !

0 m 0

= = = m, (4.262)

0 0 m

m 1 = m 1 m = I . In this case, Eq. (4.255) reduces to

D a = !2 a + i ! a (4.263)

= a. (4.264)

where D = m 1k and = !2 + i !. In this case D is symmetric because the two masses are

equal, so we can proceed using D = m 1k . Because / m the equations of motion can be

written in the standard form of an eigenvalue equation. Moreover, because m and k are the

same as for the undamped case, D is unchanged and the eigenvalue equation is formally the

same as it was for the undamped case. Thus, the eigenvalues are the same, namely = !2p

and = !2p + 2!2b . The only dierence is that for the damped case = !2 + i ! whereas

for the undamped case = !2 . Therefore the frequencies are given by solutions to

!2 + i ! !2p = 0 (4.265)

!2 + i ! !2p + 2!2s = 0 . (4.266)

From the quadratic formula we obtain

r 2

i

! = !2p (4.267)

2 2

r 2

i

! = !2p + 2!2s . (4.268)

2 2

These solutions are the same as those obtained in 2.1 for a single damped oscillator where

the normal mode frequency for each mode replaces the natural frequency !0 of the pendu-

lum. Each mode of the coupled pendulums, therefore, exhibits the same kind of behavior

as the simple pendulum, with the same underdamped, overdamped, and critically damped

cases. The only dierence is that there are two natural frequencies for the coupled pen-

dulum, namely ! and ! , whereas there is only one for the simple pendulum, namely

!0 . This leads to somewhat more complex damping for the coupled pendulums. Although

there are two normal frequencies, Eqs. (4.267) and (4.268) have only a single damping rate

/2, which is common to both modes. This is a consequence of our assumption the / m

and the fact that we took both pendulums to have the same mass.

Finally we note that the normal modes for the damped pendulum are exactly the same as

those for the undamped pendulum: there is a low-frequency mode where the pendulums

t

118 Normal Modes

oscillate in phase and a high-frequency mode where the pendulums oscillate out of phase.

The dierence in this case is simply that the oscillations die out with time.

Coupled oscillators are often driven by some external force, just as simple oscillators are,

and not surprisingly, this leads to the phenomena of resonance, with some interesting new

twists. Intuitively, we might expect that there will be a resonance at each normal frequency,

and indeed, this is what happens, albeit with a few caveats. Here we embark on a brief

treatment of forced oscillations and resonance in coupled oscillators.

To examine forced oscillations in coupled oscillators, we return once again to the prob-

lem of two coupled pendulums. We also include damping, but in such a way that the damp-

ing matrix is proportional to the mass matrix, just as we did in the previous section. This

keeps the analysis fairly simple and illustrates the basic phenomena. We will also examine

the limit of negligible damping.

To start, we include sinusoidal forcing only for pendulum 1, the left pendulum, in Fig. 4.1.

For this case, we assume no forcing of the second pendulum. Starting from Eqs. (4.260)

and (4.261), the equations of motion become

mg

m x1 = + k x1 + kx2 x1 + F0 e i!t (4.269)

l

mg

m x2 = kx1 + k x2 x2 , (4.270)

l

where F0 e i!t is a sinusoidal forcing term.

Before solving the problem, its useful to think about the underlying physics. We know

that there are two normal modes from our analysis of the undamped and damped problem:

the slow mode with the normal coordinate q = p12 (x1 + x2 ), and the fast mode with

the normal coordinate q = p12 (x2 x1 ), which are just the normalized versions of Eqs.

(4.22) and (4.23), respectively. This means that the position x1 of the driven pendulum is a

linear combination of q and q : that is, x1 = p12 (q + q ). Thus, when we drive pendulum

1 with an external force, we are exciting both q and q : that is, both the slow mode

and the fast mode. For a given drive frequency, we would not expect the two modes to

respond equally, however, as their respective responses should depend on how close the

driving frequency ! is to the normal frequencies ! and ! of the two modes. The closer

the drive frequency is to a given normal frequency, the more that mode would be expected

to respond. With these ideas in mind, lets turn to the solution of the problem.

The general approach to finding the mathematical solution to such a problem is to first

t

119 Forced oscillations

find the normal modes for the problem without external driving using the matrix formula-

tion. Thus, we start by restating the problem in matrix form. Following our development

in 4.2.1 and 4.4, the equation of motion is

i!t

m x = k x x + F0 e . (4.271)

where m and k are given by Eq. (4.181), and is given by Eq. (4.262). The column matrix

(or vector) F0 is given by

!

F0

F0 = , (4.272)

0

where the top element is finite and the bottom element is zero because only pendulum 1 is

driven.

To solve this problem, we return to the formalism developed in 4.2.2. We will also de-

velop the formalism allowing for the possibility of unequal masses as developed in 4.2.5.

The first step is to transform to a basis in which the dynamical matrix is diagonal. Following

the procedure of 4.2.5 of left multiplying both sides by m 1/2 and inserting m 1/2 m 1/2 = I

in strategic places, Eq. (4.271) becomes

m 1/2m x = m 1/2k m 1/2 m 1/2 x m 1/2 m 1/2 m 1/2 x + m 1/2 F0 e i!t . (4.273)

which reduces to

1/2 1/2 1/2 i!t

y = D y m mm y + m F0 e (4.274)

1/2 i!t

= Dy y + m F0 e , (4.275)

y = U y , (4.276)

where U T is the matrix whose columns are the normalized eigenvectors of D . Multiplying

Eq. (4.275) by U

U y = UDU T U y U y + Um 1/2

F0 e i!t

(4.277)

y = D

Dy y + Um 1/2

F0 e i!t

, (4.278)

where D 1

= U T . If there is no

damping or forcing, Eq. (4.278) becomes

y = D

Dy . (4.279)

We saw already in Eq. (??) that in the tilde basis, the equations of motion are decoupled.

This means that the matrix DD is diagonal with diagonal entries that are the eigenvalues

of the system, in this case !21 and !22 . As the equations are decoupled, we write out Eq.

(4.278) in component form,

(4.280)

y 2 + y 2 + !22 y2 = F02 e i!t

, (4.281)

t

120 Normal Modes

where

1/2

F0 = Um F0 . (4.282)

These are the same as the equation of motion for a single forced oscillator that we encoun-

tered in 3.1.1 (see Eq. (3.2)).

The eigenvectors for equal-mass coupled pendulums are

! ! ! !

sin 1 1 sin 1 1

e1 = = p , e2 = = p , (4.283)

cos 2 1 cos 2 1

where = 45 . Since the columns of U T are the eigenvectors of the system

! p ! !

1 1 1 1/ m 0 F0

F0 = Um 1/2 F0 = p p (4.284)

2 1 1 0 1/ m 0

!

1 F0 1

= p p , (4.285)

2 m 1

p

which means that F01 = F02 = F0 / 2m.

For the example here where the masses are equal, we can trivially transform between

the

Because the form of the dierential equations describing the normal modes, Eqs. (4.280)

and (4.281), is the same as the dierential equation describing a single forced oscillator,

Eq. (3.4), we can adapt the solution to Eq. (3.4) directly to Eqs. (4.280) and (4.281). Thus,

using the solution expressed in Eq. (3.15) as a guide together with Eqs. (3.12) and (3.13),

the solutions for the normal modes q and q are given by

q (t) = a (!) cos(!t ), (4.287)

where

p1 F 0 /m !

2 1 !

a (!) = p , = tan , (4.288)

(!2 !2 )2 + 2 !2 !2 !2

p1 F 0 /m

0 1

2 1B

BB ! CCCC

a (!) = q , B

= tan B@ 2 C . (4.289)

(!2 !2 )2 + 2 !2 ! !2 A

The responses of the normal modes given by Eqs. (4.286)-(4.289) are simple enough as

they mimic exactly the results we obtained for a single resonator. It is interesting to look at

the response of the individual masses, however, as they are a little less evident. For the case

of the coupled pendulums considered here, Eqs. (??) and (??) can be inverted to express

the coordinates of the individual masses in terms of the normal coordinates

1

x1 = p (q q ) (4.290)

2

1

x2 = p (q + q ) . (4.291)

2

t

121 Summary of important points of Chapter 4

4

3

2

1

A1

0

1 2 3 4

1

2

3

0

1 2 3 4

1

2

A2

t

3

4

Fig. 4.14 Resonance curves for two weakly coupled pendulums with forcing: Amplitudes of

masses 1 and 2 for the slow and fast modes as a function of forcing frequency.

These two equations together with Eqs. (4.286)-(4.289) give the steady state response of

the two masses, which are plotted in Fig. 4.14.

The motion of a system of coupled linear oscillators can be described in terms of its

normal modes.

The number of normal modes N of a system is equal to the number of degrees of free-

dom.

For each normal mode, there is a normal coordinate q that is a linear combination of

the original coordinates used to describe the system.

When the coordinates of all of the particles in the system are expressed in terms of the

normal coordinates, the system of N second order dierential equations decouples into

N uncoupled dierential equations of the form q + !2 q = 0, where ! is the (normal)

frequency of that mode.

For each normal mode there is one normal frequency.

Two or more normal modes are said to be degenerate if they have the same normal

frequency. Degenerate normal modes generally are the result of some kind of symmetry

of the system allowing the system to oscillate in essentially the same way along two or

more equivalent directions. For example, the transverse modes of the CO2 molecule are

degenerate because oscillations along the transverse directions (the y and z directions in

Fig. 4.13)) are identical.

t

122 Normal Modes

Problems

4.1 Two simple pendulums of length 0.300 m and and mass 0.950 kg are coupled by

attaching a light horizontal spring of spring constant k = 1.50 N/m to the masses.

1 Determine the frequencies of the two normal modes.

2 One of the pendulums is held at a small distance away from its equilibrium po-

sition while the other pendulum is held at its equilibrium position. The two pen-

dulums are then released simultaneously. Show that after a time of approximately

12 s the amplitude of the first pendulum will become equal to zero momentarily.

(assume g = 9.81 m/s2 ).

4.2 (French 5.5) Two identical undamped oscillators, A and B, each of mass m and

natural (angular) frequency !0 , are coupled in such a way that the coupling force

exerted on A by B is m(d2 xB /dt2 ), and the coupling force exerted on B by A is

m(d2 xA /dt2 ), where is the coupling constant of magnitude less that 1. Describe

the normal modes (i.e., specify the relative amplitudes of the displacements of mass

!20

A and mass B) of the coupled system and find their frequencies. [Ans: !2 = (1)

with respective amplitude rations AB /AA = 1.]

4.3 (French 5-6) Two equal masses on an eectively frictionless horizontal air track are

held between rigid supports by three identical springs, as shown in the figure below.

The displacements from equilibrium along the line of the springs are described by the

coordinates xA and xB , as shown. If either of the masses is clamped in its equilibrium

position, the period T = 2/! for one complete vibration of the other mass (attached

to two springs) is 3 sec.

1 If both masses are free (but still attached to the springs), what are the periods of

the two normalp modes ofpthe system? Sketch a graphs of xA (t) and xB (t) in each

mode. [Ans: 6 s and 3 2 s.]

2 At t = 0, mass A is at its normal resting position and mass B is pulled aside a

distance of 5 cm. The masses are released from rest at this instant.

1 Write an equation for the subsequent displacement for each mass as a function

of time.

2 What is the period (in seconds) for the transfer of from B to A and back. After

one pcycle is the situation at time t = 0 exactly reproduced? Explain. [Ans:

p3 ( 3 + 1)]

2

k k k

m m

4.4 Two masses m are suspended in series from springs with the same force constant k

as shown in the figure below:

t

123 Problems

(a) Find the normal frequencies and normal coordinates of this system. Hint: Define

your coordinates for the two masses with respect to the static equilibrium posi-

tions of each mass. In this way, you should find that constant terms proportional

to mg disappear from the equations of motion. [Ans:

r

1 p 1 p

!1 = (3 5) !a e1 = q (2, 1 + 5)

2 p

10 + 2 5

r

1 p 1 p

!2 = (3 + 5) !1 e2 = q ( (1 + 5), 2),

2 p

10 + 2 5

p

where !a = k/m.]

(b) At t = 0, both masses are at their respective equilibrium positions each mov-

ing with the same velocity 30 . Find expressions for the motion of each mass as

function of time given these initial conditions. Ans:

2 3

666 777

30 /!a 6

666 2 4 777

x1 (t) = p p 66 q sin ! 1 t + q sin ! t

2 7 777

2(5 + 5) 46 3 p5 3+ 5

p 5

2 3

666 p p 777

30 /!a 666 4(2 + 5)

6 2( 5 1) 777

x2 (t) = p p 66 q sin ! 1 t q sin ! t

2 7 777

2(5 + 5) 46 3 p5 3+ 5

p 5

4.5 Consider a pair of pendulums, one twice as long as the other coupled to each by a

spring. Each pendulum consists of a massless sti rod with a mass m at its end. The

spring is attached halfway along the length of the longer pendulum and to the end of

the other pendulum as shown in the figure below. Assume that 1 = 2 = 0 and that

the spring is unstretched at equilibrium.

(a) Write down the coupled equations of motion for the two pendulums, using the

angles 1 and 2 as the dynamical variables and assuming that 1 1 and 2

1. Write the coupled equations of motion in angular form using torque rF,

where I is the moment of inertia. Show that the

that is, in the form = I ,

equations of motion can be written in matrix form

m = k , (4.292)

t

124 Normal Modes

where

! ! !

ml2 0 mgl + kl2 kl2 1

m= , k= , = . (4.293)

0 4ml2 kl2 kl2 + 2mgl 2

(b) Show that the equation of motion can be rewritten in the standard form = D 0

with

!

!2 + !2 !2s

D 0 = m 1k = p 2 2 s , (4.294)

! s !2p + 2 !2s

where !2p = g/l and !2s = k/m. What is the numerical value of ?

(c) Find the symmetric form of the dynamical matrix given by D = m 1/2k m 1/2 .

p

(d) For the remaining parts of this problem set ! p = 1 and ! s = 2. Find the

normal frequencies of the system.

(e) Find the normalized eigenvectors of D . Make a plot of the eigenvectors in the

y basis defined by y = m 1/2 . Writing the normalized eigenvectors as in Eq.

(4.101), show that the eigenvectors are rotated by an angle of 67.5 in the y

basis.

(f) Find equations for the trajectories of each pendulum as a function of time for the

initial conditions where 1 = 2 = 0, 1 = > 0, and 2 = 0. Make plots of 1 (t)

and 2 (t).

4.6 In this problem you analyze the normal modes of the Wilberforce pendulum, which

consists of a body of mass m and moment of inertia I hanging from a spring, as shown

in Fig. 4.15. When a Wilberforce pendulum is set into motion, say by stretching

the spring a small distance from its equilibrium position, the pendulum begins to

oscillate up and down. As time passes, the pendant mass slowly begins to twist about

its vertical axis while the amplitude of the vertical oscillations becomes smaller due

to a weak coupling between extension and torsion of the spring. As time continues to

pass the system periodically exchanges energy between the extensional and torsional

motion.

The coupling between extension and torsion can be understood as follows. When

the mass on the spring is displaced vertically from its equilibrium position, the radius

of the coil of the spring changes slightly, causing a small bit of torsion in the spring

and changing the equilibrium twist of the spring. Thus, the potential energy of the

t

125 Problems

spring is given by

1 2 1

U(z, ') = kz + [' '0 (z)]2 , (4.295)

2 2

where '0 (z) gives the angular shift in the equilibrium twist when the spring is dis-

placed a vertical distance z from its equilibrium position. Because the displacements

are small, we can take the shift to be linear in z

1 1

U(z, ') = (k + b2 )z2 + '2 bz' . (4.297)

2 2

Typically, the coupling between extensional and torsional motion is weak, meaning

that b2 k.

(a) The vertical force exerted by the pendulum on the mass is given by Fz = @U/@z;

the torque (torsional force) is given by ' = @U/@'. Use these equations and

the expression for the potential energy above to obain the following coupled set

of equations of motion for the Wilberforce pendulum:

d2 z

m = Kz + ' (4.298)

dt2

d2 '

I 2 = z c' , (4.299)

dt

where the constants K, , and c are various combinations of k, , and b that you

should determine.

(b) What are the natural frequencies !z and !' of the translational and rotational

oscillations, respectively, in the absence of the coupling terms in Eqs. (4.298)

and (4.299), that is, when b = 0.

t

Fig. 4.15 Wilberforce pendulum twisted an angle ' '0 (z) past its equilibrium position '0 (z) for

a vertical displacement of z.

t

126 Normal Modes

s

K K 2 2

2

! = + + (4.300)

2m 2I 2m 2I mI

s

K K 2 2

!2 = + + + (4.301)

2m 2I 2m 2I mI

(d) In our class demonstration, we observed the periodic exchange of energy be-

tween vertical and torsional motion in the Wilberforce pendulum, similar to that

plotted in Fig. 4.4. One can show that beats occur like those shown in Fig. 4.4

when !z = !' . The vertical oscillation period was about 0.8 s and the beat

periodthat is the time it took the pendulum to go through one full cycle of no

torsional oscillations through a maximum in torsional oscillations and back to no

torsional oscillationswas about 20 s. Given that the mass was a steel cylinder

about 2 cm high and about 4 cm in diameter, find the approximate values of the

mass m, the moment of inertia I, the spring constant k, and the torsional spring

constant . The value of b is measured to be about 2.5 rad/m.

4.7 Consider a system exactly like the the pendulum swinging from a sliding mass pre-

sented in ?? but without a spring attached to the mass M. The mass M is still free

to slide without friction on the surface supporting it. Find the normal modes and

frequencies of the system and describe the motion of each mode physically.

4.8 Consider three equally spaced masses, M, m, and m, connected by springs and con-

strained to move on a circle of radius R that is fixed in space, as shown in the figure

below. All three springs have the same spring constant k. The dashed spokes show

particle positions with 1 = 2 = 3 = 0 corresponding to unstretched springs; dotted

lines show positive (clockwise) displacements.

(a) Write down the equations of motion and cast the problem in terms of matrices

following the formalism introduced in 4.2.1. Using the matrix formalism, find

all the normal modes of the system including the normal frequencies and the nor-

malized eigenvectors (normal coordinates). Make a sketch of the normal modes

and briefly discuss the motion associated with each one.

t

127 Problems

(b) Starting from your solution to part (a), find the equations for the motion of the

three masses for the following set of initial conditions: 1 = 0, 1 = , 2 = 0,

2 = , 3 = 0, 3 = 0,

4.9 In 4.3.1 we found three longitudinal modes of the carbon dioxide molecule, one of

which was the zero-frequency mode corresponding to simple translation of the entire

molecule at constant velocity. Reformulate the problem as one with only two degrees

of freedom by fixing the center of mass to be zero. That is, find the center of mass in

terms of x1 , x2 , x3 , m, and M. By insisting that the center of mass remain fixed, say

at a value of zero, eliminate one of the three variables x1 , x2 , or x3 , and reformulate

the problem with only two degrees of freedom. Then find the normal frequencies and

eigenvectors using the matrix formulation.

4.10 Consider the two masses hanging from springs that you considered in Problem 4.4.

Suppose the bottom mass if driven with the sinusoidal force F(t) = F0 cos !t and

that the motion of each mass is damped with a damping force Fd = 3, where 3 is

the velocity of the particular mass.

(a) Using the results from your solution to Problem 4.4 (or the one posted by your

TA), find the steady state response of the system to the external drive. In particu-

lar, find equations for the motion of the system in terms of the normal coordinates

and for the motion of the individual masses.

(b) What does the motion of the two masses look like when the drive frequency is

equal to that of the slow normal mode? the fast normal mode?

4.11 (a) In 4.1.1, we found the normal coordinates for the two pendulums coupled by

a spring. Using the matrix formalism developed in 4.2.2 show that the ma-

trix U 1 that transforms from the (x1 , x2 ) coordinates to the ( x1 , x2 ) normal co-

ordinates is a rotation through 45 . That is, show that U 1 is the conventional

two-dimensional rotation matrix (for a specific angle). For information about the

rotation matrix, see http://en.wikipedia.org/wiki/Rotation_matrix.

(b) In Appendix B, we solve for the normal modes of a double pendulum using the

Lagrangian formalism to obtain the equation of motion. Using the results of that

analysis, find the transformation matrix U 1 that transforms from the (x1 , x2 ) co-

ordinates to the (q , q ) for the double pendulum. Can that transformation be

described as a rotation between the the (, ) coordinates and the ( x1 , x2 ) normal

coordinates? That is, is the transformation matrix U 1 a rotation matrix for some

angle? Using the dot product, find the angle between ( x1 , x2 ) for the double pen-

dulum. Finally, find the normal coordinates of the double pendulum and show

by substituting them into Eqs. (B.35) and (B.36), the small-angle equation of

motion, that the equations decouple.

5 Waves of oscillating particles

Most of the waves we encounter in everyday life involve the excitation of some medium:

waves on the strings of a guitar or violin, the undulations of the surface of a pool of water

or of a drum, or sound waves in air, water, or through a solid wall. Such waves involve

particles moving up and down or back and forth, and are properly viewed as normal modes

such as those introduced in Chapter 4. There, we studied the normal modes of systems

consisting of a few particles, two, or at most three. Here, we consider waves involving a

much larger number of particles. As we shall see, the fact that particles are separated a

finite distance from each other imposes a lower cuto on the length scale on which waves

can exist, which has important physical consequences. When particles are neatly arranged

on a lattice, this leads to phenomena such as band gaps, that is, the absence of propagating

waves over certain frequency ranges. When the arrangement of particles is disordered, or

the interactions between particles vary randomly in strength, the normal modes can become

highly localized in space, involving only a subset of nearby particles within the system. In

fact, the phenomena involving waves on arrays of interacting particles is almost limitless,

and are sufficiently subtle and important to remain the focus of a great deal of current

research.

In this chapter, we begin our inquiry into the phenomena of waves of oscillating particles

by considering a deceptively simple but revealing example: point masses on a massless

string. We begin with N identical masses evenly spaced along a string, and then explore

the consequences of having N pairs of two dierent masses evenly spaced along a string,

which leads to the concept of a band gap. We then explore the consequences of other

arrangements of particles of dierent masses along a string, in particular, quasicrystalline

and disordered arrangements, both topics of current research. All of these examples are

drawn from systems that are essentially one dimensional, meaning that the particles are

always arranged along a line. This assumption simplifies the math and illustrates much of

the essential physics.

We start by considering the problem of N masses tethered to each other by a string stretched

between two fixed walls, as shown in Fig. 5.1. The masses are evenly distributed along the

string, with each mass a distance l from its nearest neighbor. At equilibrium, the masses and

string form a straight line. We shall consider the oscillations that occur when the masses

128

t

129 Identical point masses on a massless string

are displaced perpendicular (transverse) to the horizontal line connecting the masses at

equilibrium, as shown in Fig. 5.1. We index the masses by the letter s. The transverse

displacement of the sth mass is given by y s . The ends of the string are fixed to walls and do

not move so that y0 = yN+1 = 0. We consider the case where the transverse displacements

are small meaning that the angle made with he horizontal by the sth string segment s 1

(see Fig. 5.1(b)). Taking the tension T along the string to be constant, which is valid for

small transverse displacements, the equation of motion the sth mass is

Fys = T sin s T sin s 1 = my s . (5.1)

For small transverse displacements, we can make the small angle approximation that sin s '

(y s+1 y s )/l and sin s 1 ' (y s y s 1 )/l, yielding the following equation of motion

T

my s = (y s 1 2y s + y s+1 ) . (5.2)

l

Equation (5.2) actually represents N coupled equations of motion, one for each particle on

the string. Thus, as the index s runs from 1 to N, we obtain:

T

my1 = ( 2y1 + y2 ) (5.3)

l

T

my2 = (y1 2y2 + y3 ) (5.4)

l

..

.

T

myN 1 = (yN 2 2yN 1 + yN ) (5.5)

l

T

myN = (yN 1 2yN ) . (5.6)

l

Each equation has three dynamical variables, those involving a particle and its two nearest

neighbors, which couples each equation to the equation before and after it. The only excep-

tions are the first and last equations, which have only two terms because of the boundary

conditions y0 = yN+1 = 0 discussed above.

We can recast Eqs. (5.3)-(5.6) as a matrix equation of the same form we encountered in

Chapter 4:

m y = k y , (5.7)

t

Fig. 5.1 N masses on a string under tension T . At equilibrium, the masses form a straight

horizontal line. (a) Transverse displacements {y s } of N particles of mass m separated

by a horizontal distance l on a string. (b) Close-up showing forces on the sth mass m

tethered to a string under tension T .

t

130 Waves of oscillating particles

where

0 1 0 1 0 1

BBBm 0 CC

C BBB1 0CC

CCC BBB y1 CCC

BBB CCC BBB BBBB y2 CCCC

BBB m CCC BBB 1 CCC BBB CCC

BBB C

C BBB CCC B C

m = BBB m C

CCC = m BBB 1 CCC , y = BBBB y3 CCCC

BBBB CCC BBB CCC BBBB .. CCCC

.. CCC BBB .. CCC

BBB . C B . CCA BBB . CCC

@ A @ @ A

0 m 0 1 yN

0 1 0 1

BBB 2T/l T/l 0 CCC BBB 2 1 0 CCC

BBB CCC BBB CCC

BBB T/l 2T/l T/l CCC BBB 1 2 1 CCC

BBB .. CCC T BBB .. CCC

k = BBB B T/l 2T/l . CCC = BBB 1 2 . CCC .

BBB CCC l BBB CCC

BBB .. .. CC BBB .. .. CC

BBB . . T/lCCCCC BBB

B@

. . 1CCCCC

@ A A

0 T/l 2T/l 0 1 2

With these definitions, performing the matrix multiplications implied by Eq. (5.7) yields

Eqs. (5.3)-(5.6). There are many techniques for solving Eq. (5.7), including those we de-

veloped in Chapter 4. However, the methods developed in Chapter 4 become increasingly

inefficient as the number of particles N grows and the sizes of the matrices grow with

them. For greatest efficiency, we need to take advantage of some special properties of the

matrices m and k , including their symmetry. We defer that discussion for now, however,

in favor of a simpler method based on the physics of the problem. We shall return to the

more general method afterwards, however, as it allows us to solve a much broader range of

problems.

We can gain some insight into the problem of N point masses on a string by considering

the small N limit. Lets start with N = 1. Because there is only one particle, there is but one

degree of freedom and one equation of motion, which we can get directly from Eq. (5.2):

2T

y1 , my1 = (5.8)

l

where we have set the positions of the ends of the string to zero: y0 = y2 = 0. This is the

p

equation for a simple harmonic oscillator with an oscillation frequency of !(1)

= 2T/ml.

Here the superscript refers to the number of particles on the string, in this case 1; the

subscript refers to the normal frequency, labeled in Greek alphabetical order. For this trivial

case, there is one degree of freedom and one normal mode.

Lets consider the case of two particles on a string: N = 2. In this case, there are two

equations of motion, one each corresponding to particles s = 1 and s = 2 in Eq. (5.2),

T

( 2y1 + y2 ) my1 = (5.9)

l

T

my2 = (y1 2y2 ) , (5.10)

l

where y0 = y3 = 0. Based on our experience with normal modes in Chapter 4, it is pretty

easy to guess what the normal modes must look like: a slow () mode where both particles

t

131 Identical point masses on a massless string

oscillate together (in phase) and a fast ( ) mode where the two particles oscillate in opposi-

tion (exactly out of phase). These two modes are depicted in Fig. 5.2. Thus we would guess

that the normalized normal coordinates are given by q = p12 (y1 + y2 ) and q = p12 (y1 y2 ).

Exercise 5.1.1 By adding and subtracting Eqs. (5.9) and (5.10), show that you can

obtain two decoupled equations of motion where the normal coordinates are given

by q = p12 (y1 + y2 ) and q = p12 (y1 y2 ). What are normal frequencies of the two

normal modes?

For the case of three particles on a string, N = 3, the equations of motion become

T

my1 = ( 2y1 + y2 ) (5.11)

l

T

my2 = (y1 2y2 + y3 ) (5.12)

l

T

my3 = (y2 2y3 ) , (5.13)

l

where y0 = y4 = 0. We already treated this problem in 4.2.4, but even if you do not recall

the exact results, you can probably guess the qualitative form of the particle oscillations

corresponding to the three normal modes. You can compare your intuition against the Fig.

5.2, which shows diagrams of the various normal modes for the cases where there are 1 to

3 particles.

Exercise 5.1.2 Write down the equations of motion for the case where there are four

equally spaced particles of mass m on a string. By extending the diagrams in Fig.

5.2, try to guess the normal coordinates for the case where there are four particles.

N =1 n=1 N =3 n=1

N =2 n=1 N =3 n=2

N =2 n=2 N =3 n=3

t

Fig. 5.2 Normal modes for N masses on a string for N = 1, 2, 3.

t

132 Waves of oscillating particles

Test your guesses by forming sums of from the linear combinations of y1 , y2 , y3 , and

y4 that form the normal coordinates. A correct guess should lead to an equation of

motion that is a function of a single normal coordinate.

Returning to the general problem of N particles on a string, we begin by assuming a trial

solution to Eq. (5.2) of the usual form:

i!t

ys = as e . (5.14)

Here y s and a s represent the instantaneous transverse displacement and the amplitude,

respectively, of the sth particle. By adopting this trial solution we are assuming, as we

always do in normal mode analysis, that all particles are oscillating at the same frequency

! for a given normal mode. Substituting Eq. (5.14) into Eq. (5.2) yields

!

T 2T T

as 1 + !2 m a s a s+1 = 0 (5.15)

l l l

p

Dividing through by m and defining the characteristic frequency !0 = T/ml, this be-

comes

!20 a s 1 + 2!20 !2 a s !20 a s+1 = 0 . (5.16)

Equation (5.16) actually represents N equations, one for each value of s, subject to the

boundary conditions that both ends are fixed at zero displacement: a0 = aN+1 = 0. While

solving this set of N coupled equations is straightforward for small values of N, it becomes

quite difficult for large values of N. Therefore, to assist us in our search for a solution, we

take a clue from the low-N solutions discussed above, which suggest a spatially oscillating

solution. So we start by assuming an oscillatory trial solution, say of the form sin s or

cos s. The algebra is made simpler if we write our trial in exponential form

a s = Aei(s )

, (5.17)

where we take A to be real and explicitly include the possibility of non-zero phase . As

usual, we are concerned only with the real part of a s . The parameter is unknown at this

point but its value clearly determines the wavelength of the oscillation. Therefore, for a

given number of particles N, we would expect to depend on the mode n.

First, lets apply the boundary conditions at s = 0 and s = N + 1 to our trial solution:

i

a0 = Ae =0 (5.18)

i[(N+1) ]

aN+1 = Ae =0. (5.19)

Recalling that we are only concerned with the real part of a0 , the boundary condition at

s = 0, Eq. (5.18), implies that = /2, or some other odd multiple of /2. Adopting the

simplest choice = /2, the boundary condition at s = N + 1 becomes

aN+1 = Aei[(N+1) /2]

= A sin[(N + 1)] = 0 , (5.20)

where once again we recall that we are only concerned with the real part of aN+1 . Equation

(5.20) is satisfied only if (N + 1) = n, where n = 1, 2, 3, ..., which means that

n

= where n = 1, 2, 3, ... (5.21)

N+1

t

133 Identical point masses on a massless string

ns

a s = A sin where n = 1, 2, 3, ... (5.22)

N+1

For each value of n, Eq. (5.22) gives the N particle positions as s runs from 1 to N. Thus,

each integer n must denote a distinct normal mode. Since there are N particles, and thus

N degrees of freedom, there are N normal modes. This means that n, which indexes the

normal modes, will be restricted to run from n = 1 to N. We will explore the allowed values

of n in greater detail below, but defer that discussion for the moment.

We still need to check to see if our trial solution satisfies the Eq. (5.16). To this end, we

substitute Eq. (5.17), which gives

!20 e i + 2!20 !2 !20 ei = 0 , (5.23)

where we have canceled out the common factors of Aei(s ) . Solving for !2 gives

!2 = 2!20 !20 ei + e i (5.24)

= 2!20 (1 cos ) (5.25)

= 4!20 sin2 , (5.26)

2

where we have used the identity 1 cos = 2 sin2 (/2). Using the expression for in Eq.

(5.21), Eq. (5.26) becomes

n

!n = 2!0 sin , (5.27)

2(N + 1)

where we have added a subscript to !n to denote that n = 1, 2, 3, ..., N, which gives N

distinct normal frequencies, one for each of the N normal modes.

Figure 5.3 shows a plot of the (non-negative) normal frequencies as a function of mode

number n for n 0 as given by Eq. (5.27). Here we include all positive values of n. There

are N + 1 distinct frequencies between n = 0 and n = N. There are no new or values of !n

outside this range; in fact, you can show from Eq. (5.27) that !n+2s(N+1) = !2s(N+1) n = !n

20

t

n

0 10 20 30 40 50 60 70 80

Fig. 5.3 Normal frequencies vs. mode number n for N = 16 identical masses on a string. Dark

circles denote modes for 1 n 16.

t

134 Waves of oscillating particles

for 0 n N + 1 where s is any integer. As we will see below, the particles do not

move for the modes n = 0 and n = N + 1. Thus, although mathematically n can take on

any integral value, there are only N distinct normal modes with finite displacements, those

corresponding to 1 n N.

The amplitude dependence given by Eq. (5.22) is for a single normal mode. Because

there are N dierent normal modes indexed by n, we add a subscript n to our equation for

the amplitude of the sth particle for the nth normal mode:

ns

ans = An sin where n = 1, 2, 3, ..., N (5.28)

N+1

Taken together, Eqs. (5.27) and (5.28) constitute the full general solution to the problem.

To obtain the motion of each of the particles in a given normal mode, we simply substi-

tute these expressions into Eq. (5.14), including both the positive and negative frequency

solutions allowed by Eq. (5.27),

ns

yns (t) = sin An ei!n t + A0n e i!n t . (5.29)

N+1

Alternatively, we can write the sum of complex exponential terms as a sum of a sine and

cosine terms

ns

yns (t) = sin (Bn sin !n t + Cn cos !n t) . (5.30)

N+1

The full solution is obtained by summing over all normal modes

N

X N

X ns

y s (t) = yns (t) = sin (Bn sin !n t + Cn cos !n t) , (5.31)

n=1 n=1

N+1

where the set of {Bn } and {Cn } are determined by the initial conditions. Equation (5.31)

constitutes a complete solution to the problem of N identical masses on string.

It is instructive to examine the normal modes for at least one value of N. Lets look

at the case of N = 5, that is, five particles of mass m evenly spaced on a string. The

instantaneous displacements as given by Eq. (5.31) are shown for n = 1 to 12 in Fig.

5.4 (with Bn = 0 and Cn = 1). First we note that modes 6 and 12 are null modes for

which there are no displacements. These represent trivial solutions (no motion) and this

are ignored. Next, notice that modes 712 are the same as modes 1-5, but in reverse order

and with a trivial 180 phase dierence. You can also check that modes with the same

pattern of displacements oscillate at the same frequency, for example, !5 = !7 ' 1.93!0 .

Thus we see that there for N = 5, there are only 5 independent modes, which means we

only need to consider mode numbers n between 1 and N (= 5 in this case).

The eigenvector of the nth mode is given by an = (an1 , an2 , ..., an5 )T for the case where

N = 5. According to Eq. (5.28) the components ans are proportional to sin[ns/(N + 1)].

P

Normalizing the eigenvectors so that aTn an = Ns=1 a2ns = 1, the eigenvectors are given by

r

2 ns

ans = sin , s = 1, 2, ..., N . (5.32)

N+1 (N + 1)

t

135 Identical point masses on a massless string

n=1 n=7

n=2 n=8

n=3 n=9

n=4 n = 10

n=5 n = 11

n=6 n = 12

t

Fig. 5.4 Normal modes for five particles on a string (N = 5) shown when each particle is at its

maximum amplitude (dark) and at different times during each cycle (light). Modes

7-11 are not new modes but merely repeat modes 1-5, but in reverse order. Modes 6

and 12 are null modes that have zero amplitude (i.e. they do not oscillate).

For N = 5 particles, the five eigenvectors {an } given by Eq. (5.32) are

0 1 1 0 1 1 0 1 1 0 1 1 0 1 1

BBB 2 p3 CCC BBB 2 CCC BBB p3 CCC BBB 2 CCC BBB 2 p3 CCC

BBB 1 CCC B

B 1 C

C B

B

B C

C

C B

B 1 C

C BBB 1 CCC

BBBB 2 CCCC BBBB 2 CCCC BBBB 0 CCCC BBBB 2 CCCC BBB 2 CCC

BB 1 CC

B 1 C B C B 1 C B C

a1 = BBBB p3 CCCC , a2 = BBBB 0 CCCC , a3 = BBBB p3 CCCC , a4 = BBBB 0 CCCC , a5 = BBBB p3 CCCC . (5.33)

BBB 1 CCC BBB 1 CCC BBB CCC BBB 1 CCC BBB 1 CCC

BBBB 2 CCCC BBB@ 2 CCCA BBBB 0 CCCC BBB@ 2 CCCA BBBB 2 CCCC

@ p1 A 1 @ p1 A 1 @ p1 A

2 3 2 3 2 2 3

t

136 Waves of oscillating particles

The entries for each eigenvector track the corresponding amplitudes for the five modes

plotted in Fig. 5.4, as expected.

Equation (5.31), where the frequencies are given by Eq. (5.27), represents the general

solution to the eigenvalue problem. To find the special solution for a specific case, one

must specify the initial conditions, meaning the initial positions and velocities of all N

particles. With this information, the constants {Bn } and {Cn } appearing in Eq. (5.31) can be

determined for any system of N particles on a string. Once the {Bn } and {Cn } are known,

Eq. (5.31) provides a complete description of the subsequent motion of the particles. In

practice, the algebra associated with doing so can be quite tedious for any system made up

of more than 3 or 4 particles, unless the initial conditions possess some special symmetry

that allows one to determine some large fraction of the coefficients by inspection. Problem

5.1 at the end of the chapter explores the use of symmetry to greatly reduce the number

of coefficients that must be determined analytically. When such symmetry is absent, the

task of determining the coefficients {Bn } and {Cn } is best left to computers! We take up this

issue, and much more, in the remainder of this chapter.

Before moving on to consider more complex problems involving masses on strings, we

note that the horizontal distance along the string of a particle with index s is x s = sl. Using

this fact, we can replace s in in Eq. (5.31) by x s /l and rewrite our solution in terms of the

particle coordinate

N

X N

X !

nx s

y s (x s , t) = yns (t) = sin (Bn sin !n t + Cn cos !n t) (5.34)

n=1 n=1

(N + 1)l

N

X

= (sin kn x s )(Bn sin !n t + Cn cos !n t) , (5.35)

n=1

where

n

kn = . (5.36)

(N + 1)l

We recognize kn as the wave vector of the nth mode, with a corresponding wavelength of

2(N + 1)l/n. More importantly, the form of Eq. (5.35) is that of standing waves that are

made up of sums of counter propagating traveling waves propertional cos(kn x !n t) and

sin(kn x !n t).

We can also express the normal frequencies from Eq. (5.27) in terms of kn

1

!n = 2!0 sin kn l . (5.37)

2

For the lowest n modes where n N, we can make the approximation sin 12 kn l 12 kn l, in

which case Eq. (5.37) becomes

!n = !0 kn l . (5.38)

t

137 Identical point masses on a massless string

Thus, for small n, which corresponds to the longest wavelength normal modes, the fre-

quency increases linearly with the wave vector kn . In 5.3, we show that this linear rela-

tionship between !n and kn is consistent with a wave moving with a velocity vn = !n /kn =

p

!0 l = T l/m. As we shall see, the fact that frequency ceases to increase linearly with kn

(and n) has important consequences for the propagation of pulses in such a system, but we

defer discussion of wave velocity and pulse propagation to 5.3.

The method used in the previous section works well when the high degree of physical

symmetry in the problem, as when the masses are identical and equally spaced, allows

you to guess overall form of the solution. When the particle masses or spring constants

vary from particle to particle, guessing a solution becomes more difficult and practically

impossible in many cases. We need a more general method for solving such problems.

The matrix formulation for solving normal mode problems introduced in Chapter 4

works. You simply write down the equations of motion in matrix form and solve them,

as outlined in 4.2. Once again we start from Eq. (4.180), which we rewrite here for con-

venience

m x = k x . (5.39)

We will limit our discussion to those cases where the mass matrix is diagonal, but will in

general allow for cases where the masses might be dierent from one another

0 1

BBBm1 0 CC

CCC

BBB

BBB m2 CCC

BBBB CCC

m = BB m3 CCCC (5.40)

BBB . CCC

BBB .. CCC

B@ CA

0 mN

i!t

Following the usual procedure, we posit a trial solution a = ae . substituting this into

Eq. (5.39), we obtain

!2 m a = k a . (5.41)

In order to obtain a symmetric dynamical matrix, we once again define m 1/2 , which again

p

is simply the diagonal matrix with entries mi along the diagonal. Following the same

procedure we developed in 4.2.5, we transform Eq. (5.39) to obtain

Db = b, (5.42)

where = !2 and

1/2 1/2

D =m km (5.43)

b = m 1/2 a . (5.44)

t

138 Waves of oscillating particles

Our task is now, as it was before, to find the eigenvalues and eignevectors of the symmetric

dynamical matrix D , which give us the normal frequencies and coordinates.

The first step towards finding the normal modes and frequencies of our eigenvalue equation

Eq. (5.42) is to consider a certain class of transformations called similarity transformations.

We start with a matrix S defined so that its inverse S 1 b = b0 , which of course means that

S b0 = b. If we substitute S b0 for b in Eq. (5.42) and left multiply both sides by S 1 , we

obtain

S 1D S b0 = S 1S b0 . (5.45)

D 0 = S 1D S , (5.46)

D 0 b0 = b0 . (5.47)

the two matrices D and D 0 have the same set of eigenvalues. Transformations of the form

S 1D S , where S is some invertible matrix, preserve eigenvalues and are called similarity

transformations. This property of similarity transformations, namely that they preserve

eigenvalues, makes them particularly useful for finding solutions to eigenvalue problems.

Suppose, for example, that we could find a matrix S such that the matrix formed by the

similarity transformation D 0 = S 1D S is a diagonal matrix. Written out explicitly, Eq.

(5.47) would then be

0 0 10 1 0 01

BBBD11 0 CC BB b01 CC BBB b1 CCC

BBBB 0

C

CCCC BBBB b0 CCCCB C BBBB b0 CCCC

BBB D 22 CCC BBB 2 CCC BBB 2 CCC

BBB 0

D33 CCC BBB b CCC 0 B 0C

BBB CCC BBB 3 CCC = BBBBB b3 CCCCC . (5.48)

BBB .. CCC BBB .. CCC BBB .. CCC

BBB . CCC BBB . CCC BBB . CCC

@ A@ A @ 0A

0 D0NN b0N bN

We can obtain the eigenvectors and eigenvalues to this problem by inspection. For b1 =

(1 0 0 ... 0)T the eigenvalue is 1 = D011 . Similarly, for b2 = (0 1 0 ... 0)T , the eigenvalue is

0

2 = D22 . Continuing this pattern, we see that the diagonal elements of D correspond the

the eiqenvalues, that is,

j = D0j j , (5.49)

with the corresponding eigenvector consisting of 1 in the jth entry and zeros everywhere

else. If we were able to find the matrix S that diagonalized D , then the eigenvectors for the

original matrix D would be given by

b = S b0 , (5.50)

t

139 Identical point masses on a massless string

where b0j is a column vector with 1 in the jth position and zeros everywhere else. In fact,

because the columns of the identity matrix I are the eigenvectors of D 0 , then according to

Eq. (5.50), the columns of S I = S are the eigenvectors of D ; that is, the columns of S

are the eigenvectors of D . Thus we see that solving the eigenvalue problem posed by Eq.

(5.42) is is equivalent to finding the matrix S that diagonalizes D through the similarity

transformation S 1D S .

For all but the smallest or simplest matrices, diagonalizing D through similarity transfor-

mations is done numerically. There are a number of dierent algorithms for accomplishing

this task. Most work by repeatedly applying similarity transformations that systematically

reduce the size of the o-diagonal elements of D until they are sufficiently close to zero

that the matrix is diagonalized to whatever precision is desired. The first step of such an

algorithm is to apply a similarity transformation

D 1 = S 1 1D 0 S 1 , (5.51)

which, for a well-chosen matrix S1 , makes the o-diagonal elements, or at least some frac-

tion of the o-diagonal elements of D 1 smaller than those of D 0 . The algorithm continues

in this fashion systematically, next with two transformations,

D 2 = S 2 1D 1 S 2 (5.52)

1 1 0

= S2 S1 D S1 S2 , (5.53)

which, on average, further reduces the magnitude of the o-diagonal matrix elements of

D 2 . This process is then repeated as many times as required to eect full diagonalization

Defining the matrix of the accumulated transformation

S = S 1 S 2 S 3 ... , (5.55)

1

whose inverse is given by S = ... S 3 1S 2 1S 1 1 , we see that S diagonalizes D 0 by a similar-

ity transformation

D diag = S 1D 0 S . (5.56)

In fact, the matrix S not only diagonalizes D 0 , the columns of S are the eigenvectors of D 0 ,

which follows directly from Eq. (5.50). Thus, the diagonalization process gives both the

eigenvalues of D 0 , which are the diagonal elements of D diag = S 1D 0 S , and the eigenvec-

tors of D 0 .

There are a number of algorithms that can be used to obtain S and diagonalize D 0 . Here we

describe the Jacobi method. It is one of the most basic and robust methods; it nearly always

works for the real symmetric matrices we are considering here. Other faster methods exist

t

140 Waves of oscillating particles

and are recommended when working with large matrices, but the Jacobi method suffices

for our purposes.1

The basic idea of the Jacobi method is to perform a series of planar rotations that suc-

cessively zeros the o-diagonal elements of D 0 one by one. One concern, however, is that a

transformation that zeros one o-diagonal matrix element can undo some of the previously

zeroed matrix elements. Nevertheless, it turns out that on average, the o-diagonal ele-

ments get smaller with each successive Jacobi rotation and go to zero to within any desired

precision, limited only by the precision of the computer.

The Jacobi method uses the N-dimensional generalization of the 2-dimensional rotation

matrix (see Appendix A), which is called the Jacobi rotation matrix:

0 1

BBB 1 CCC

BBB .. CCC

BBB . CCC

BBB CCC

BBB 1 CCC

BBB CCC

BBB c s CCC

BBB CCC

BBB 1 CCC

BBB CCC

.. ..

R = BBBB .

..

. .

CCC

CCC (5.57)

BBB CCC

BBB 1 CCC

BBB CCC

BBB s c CCC

BBB CCC

BBB 1 CCC

BBB CCC

BBB .. CCC

BBB . CA

@

1

where s = sin , c = cos , and is the rotation angle (explained below). The matrix R

has the same size as m and k . All the diagonal elements of R are 1 and all the o-diagonal

elements are 0 except in two rows and columns, which we label i and j (we take i < j): the

diagonal elements Rii = R j j = c and the o-diagonal elements Ri j = R ji = s. The Jacobi

rotation matrix is orthogonal, meaning R 1 = R T , making the similarity transformation

A 0 R 1AR (5.58)

T

= R AR (5.59)

changed. In fact, the entire matrix A 0 is unchanged except for the entries in rows i and j and

columns i and j. Performing the matrix multiplications, one finds that the matrix elements

1 The approach described here follows Numerical Recipes by Press et al. While significant speed gains can be

realized using more advanced methods for matrices as small as 10 10, one can readily diagonalize matrices

as large as 32 32 using the approach outlined here without unduly taxing the patience of most users. See

Numerical Recipes by Press et al. for more information.

t

141 Identical point masses on a massless string

A0ii = c2 Aii + s2 A j j 2scAi j (5.60a)

2 2

A0j j = s Aii + c A j j + 2scAi j (5.60b)

2 2

A0i j = A0ji = (c s )Ai j + sc(Aii A j j) (5.60c)

A0ik = A0ki = cAik sA jk , k,i, k, j (5.60d)

A0jk = A0k j = cA jk + sAik , k,i, k, j (5.60e)

In the Jacobi method we choose the rotation angle so that the o-diagonal matrix element

A0i j = 0. From Eq. (5.60c) this leads to the following condition

A j j Aii c2 s2 c s 1

= = = t, (5.61)

Ai j sc s c t

where t = s/c = tan . Defining the ratio of matrix elements as

A j j Aii

2 = , (5.62)

Ai j

Eq. (5.61) becomes

t2 + 2t 1=0. (5.63)

Either root of Eq. (5.61) can be used in our matrix diagonalization algorithm but the smaller

(in absolute value) of the two roots, which corresponds to a rotation angle || < 45 , gives

the more numerically stable solution. This can be written as

sgn ()

t= p . (5.64)

|| + 2 + 1

For large values of , some computational cost can be saved by replacing Eq. (5.64) by

1

t= , (5.65)

2

which has the additional benefit of avoiding a potential overflow in the computation of 2 .

Once t is known, c and s can be obtained from trigonometric identies, which give

1

c= p , s = ct. (5.66)

1 + t2

Exercise 5.1.4 Roundp o error: The more familiar way to write the solution to Eq.

(5.64) is t = 2 + 1. Convince yourself that we can write the smaller of these

two roots as

p

t = sgn () || + 2 + 1 . (5.67)

Show that Eq. (5.64) is equivalent to Eq. (5.67). Then write a computer program to

plot t vs. on a log-log plot for ranging form 10 1 to 1010 . You should find that

Eq. (5.67) yields nonsensical numerical values of t for sufficiently large values of ,

while Eq. (5.64) yields accurate values. The reason is that Eq. (5.67) is susceptible

to round-o error, which stems from the finite precision with which numbers are

t

142 Waves of oscillating particles

for Eq. (5.67) but not for Eq. (5.64). What lesson do you learn from this example

about minimizing round-o error?

For computation of A 0 , the formulas in Eq. (5.60) can be simplified, leading to more

efficient computation and less round-o error. Starting with Eq. (5.60c), we have chosen t

so that the o-diagonal element A0i j = 0, so we replace Eq. (5.60c) by

A0i j = 0 . (5.68)

Setting A0i j = 0 in Eq. (5.60c) and using the result to eliminate A j j from Eq. (5.60a) gives

where we have used the trigonometric identity s2 + c2 = 1. Equation (5.69) expresses the

updated value of A0ii in terms of the old value Aii plus a small correction, since the o-

diagonal matrix element Ai j tends towards zero for successive iterations of the algorithm,

and thus promotes stable convergence of A0ii . Continuing this approach,

A0ik = A0ki = Aik s(A jk + Aik ) , k,i, k, j (5.71)

A0jk = A0k j = A jk + s(Aik A jk ) , k,i, k, j, (5.72)

where

s

= . (5.73)

1+c

A single application of a Jacobi rotation will, as advertised, zero one o-diagonal el-

ement Ai j . Therefore it is applied successively to dierent o-diagonal matrix elements.

As noted above, successive applications can undo some of the previously zeroed matrix

elements but the net eect is that the matrix elements get smaller and tend towards zero.

There are dierent strategies for how to implement the Jacobi rotation method. Jacobis

original idea was to zero the largest o-diagonal elements first. While searching through

the o-diagonal elements for the largest one at each step is a not the most efficient strategy,

we adopt it here because of its computational simplicity. Faster schemes exist, including

those that take advantage of other special symmetries of the matrix D appearing in the

eigenvalue equation Eq. (??). For example, for many of the problems we consider in this

chapter, D is tri-diagonal, meaning that all the matrix elements are 0 except the diagonal

elements and those one o the diagonal in either direction. Very efficient schemes exist for

solving the eigenvalue problem for such matrices. If you are interested in exploring these

schemes, a good place to start is Numerical Recipies by Press et al.

The Jacobi rotation method is successively applied to the whatever the largest matrix

element is in the entire upper right (or lower left) of the matrix D after the latest rotation

until all the matrix elements are smaller than some pre-set tolerance, typically some very

small fraction (e.g. 10 9 ) of the average diagonal element. At this point, the matrix is

eectively diagonalized and the diagonal elements are the eigenvalues to within the pre-set

tolerance.

t

143 Identical point masses on a massless string

The eigenvectors are given by the columns of the matrix S , which according to Eq.

(5.55), is the product of the transformation matrices that diagonalize D . For the Jacobi

method, transformation matrices are just the Jacobi rotation matrices. We can keep track of

the cumulative product by initializing S to be the identity matrix and calculating S 0 = SR

at each step of the diagonalization procedure, where S is the cumulative product after the

previous rotation. For each rotation, only columns i and j are changed. The changes are

given by

B0k j = Bk j s(Pki + Pk j ) . (5.75)

The Jacobi diagonalization procedure gives all the eigenvalues and eigenvectors, albeit

in no particular order. Therefore, one usually takes the additional step of sorting the eigen-

values and eigenvectors from smallest to largest eigenvalue. Finally, as the eigenvectors are

only determined to within some scale factor, it is useful to normalize the eigenvectors.

Lets review the procedure developed for finding the eigenvalues and eigenvectors:

Use Choleski decomposition to find L such that m = LL T . [Appendix A]

Form D = L 1k (L 1 )T . [Eq. (??)]

1

Need to find L to form D . [See Appendix A for matrix inversion routines.]

LT

Save (L ) 1

for later use [a = (L T ) 1 a].

Use the Jacobi method to find solutions to D a = a.

Perform Jacobi rotation on the largest o-diagonal element until all the o-diagonal

elements are smaller than some pre-set tolerance.

For each rotation, update the transformation matrix using S 0 = SR .

When the Jacobi method converges, extract the eigenvalues { } and eigenvectors

{a} of D .

Eigenvalues { } are diagonal elements of D .

Eigenvectors {a} are columns of S .

Convert eigenvectors a of D to eigenvectors of m and k using a = (L T ) 1 a. [Eq. (??)]

Sort eigenvalues and eigenvectors in order from lowest to highest eigenvalues {!}.

Normalize eigenvectors.

Get normal mode frequencies from eigenvalues using { } = {!2 }.

Appendix C. As a test of the procedure and the code implementing it, we present in Fig.

5.5 the results of using this numerical procedure to solve for the normal modes of a string

under tension with 32 equally spaced identical masses compared to the analytical solution

we developed in 5.1.1. As you can see, the numerical solution reproduces the analytical

solution quite well, at least to within the accuracy that is discernible from the plots in Fig.

t

144 Waves of oscillating particles

5.5. In fact, the solutions are accurate to virtually any desired precision, certainly to the 15

or 16 significant figures typical of computer computations.

In the next sections, we use the computer code used above, and modifications to it, to

explore the normal modes that can arise for dierent distributions of masses on a string

held under tension. As we shall see, a host of new and interesting phenomena emerge from

this project.

Thus far we have considered only the simplest case of masses on a string: identical masses

equally spaced along the string. A variety of interesting new phenomena arise when we

allow the masses along the string to have dierent values. We shall consider three dierent

cases: (1) there are two dierent masses, m1 and m2 , that are arranged in an alternating

periodic sequence, (2) there are two dierent masses, m1 and m2 , that are arranged in a

Fibonacci sequence, which corresponds to a quasi-periodic sequence, and (3) a random

sequence of masses.

l apart on a string under tension T .

Numerical solution

The problem is set up exactly as in 5.1 except that the mass matrix m has an alternating

sequence of masses along its diagonal

0 1

BBB m1 CCC

BBB CCC

BBB m 2 0 CCC

BBB CCC

BBB m1 CCC

m = BBB . CCCC . (5.76)

BBB . . CCC

BBB CCC

BBB 0 m1 CCC

@B A

m2

It is simple enough to incorporate this change into our computer code and to solve for the

set of normal frequencies {!2 }. The result is shown in Fig. 5.6, where we have used the

following parameter values: m1 = 10 g, m2 = 2.5 g, l = 4 cm, and T = 60 N. The most

striking feature of Fig. 5.6 is the absence of normal frequencies between 548 and 1095

rad/s. Instead we see that there are two bands of normal frequencies, one from zero to

about 548 rad/s and another from 1095 to 1200 rad/s. The range of frequencies from 548

and 1095 rad/s is known as a band gap. There are no normal modes with frequencies in the

range of the band gap.

t

145 Sequences of different point masses on a string

ys

0.4 1

0.2

0.0 s

8 16 24 32

0.2

ys

0.4 17

0.2

0.0 s

8 16 24 32

0.2

0.4

ys

0.4 32

0.2

0.0 s

8 16 24 32

0.2

0.4

1200

1000

800

600

400

200

0 n

t

8 16 24 32

Fig. 5.5 Normal modes for a system of 32 identical particles on a string. Top three graphs

show modes 1, 17, and 32, respectively. The dotted lines visible for modes 17 and 32

show the mode amplitudes given by Eq. (5.32). Bottom plot shows frequencies vs.

mode number.

To develop some physical understanding of what is happening near the band gap, we

examine the lowest-n normal modes, modes 1 and 2, as well as the modes just below and

just above the band gap, that is, modes 16 and 17, respectively, which are shown in Fig.

5.7. The n = 1 and n = 2 modes are the longest wavelength modes, and are just what you

might expect. Their forms are hardly aected by the dierent masses on the string: simple

sinusoidal waves of wavelength 2L and L, respectively, where L = (N + 1)l is the distance

between the fixed ends of the string.

t

146 Waves of oscillating particles

As the mode number n increases, the wavelength of the modes becomes progressively

shorter, until by mode 16, the wavelength is essentially 4l, twice the distance between two

small or two large particles. There are two ways to create a wave of wavelength 4l: either

the large particles undergo large amplitude oscillations with the small particles undergoing

very small amplitude oscillations near the nodes, or the small particles undergo large am-

plitude oscillations with the small particles undergoing very small amplitude oscillations

near the nodes. These two modes of vibration corresponds to normal modes 16 and 17 re-

spectively, as shown in Fig. 5.7. Given this observation, it is straightforward to estimate the

normal frequencies of these two modes: both modes are governed approximately by Eq.

(5.8), the equation of motion for a single particle on a string oscillating between two fixed

points (see Fig. 5.2 diagram for N = 1): for mode 16, a heavy particle of mass m2 oscillates

between two between two light particles of mass m1 that are nearly stationary, while for

mode 17, it is is the reverse. Applying Eq. (5.8) to these two cases, the frequencies of the

dierent modes are given approximately by

r

2T

!16 ' = 548 s 1 (5.77)

m2 l

r

2T

!17 ' = 1095 s 1 , (5.78)

m1 l

which, to three significant figures, are the same results obtained from the complete normal

mode analysis. It is this step change in the frequency, associated with the disparity in

masses m1 and m2 but without a corresponding change in the wavelength, that leads to the

band gap.

Above the band gap, that is for modes 17 through 32, the frequency is relatively flat.

To see why, lets look at the highest n modes, which we show in Fig. 5.8. There we see

that modes 30-32, the three highest n modes, all have a basic wavelength of 2l, but that

(s 1)

1200

1000

600

400

200

t

0 n

0 5 10 15 20 25 30

Fig. 5.6 Frequencies of normal modes vs. mode number for alternating sequence of

evenly-spaced masses m1 and m2 on a string.

t

147 Sequences of different point masses on a string

ys

0.4 1

0.2

0.0 s

8 16 24 32

0.2

ys

0.4 2

0.2

0.0 s

8 16 24 32

0.2

0.4

ys

0.4 16

0.2

0.0 s

8 16 24 32

0.2

0.4

ys

0.4 17

0.2

0.0 s

8 16 24 32

0.2

0.4

t

Fig. 5.7 Normal modes 1, 2, 16, and 17, indicated on upper right, of of 32 alternating unequal

masses on a string.

the overall modulation of the amplitudes is dierent for each mode. In fact, if we consider

each particle type separately, we see that the wavelength of the small particles only is 2L for

mode 32, L for mode 31, and 2L/3 for mode 30. Apparently, the eect of the modulations

is to decrease the oscillation frequency, but only by a small amount.

Analytical solution

In the numerical solutions obtained above, we found that for a given normal mode, the

alternating masses m1 and m2 each oscillate with the same periodicity (i.e. wavelength).

For the first 16 modes, the masses m1 and m2 oscillate in phase with each other, while

for the second 16 modes, they oscillate out of phase with each other, and with dierent

amplitudes. This suggests that we may be able to find analytic solutions to the equations

t

148 Waves of oscillating particles

ys

0.4 32

0.2

0.0 s

8 16 24 32

0.2

0.4

ys

0.4 31

0.2

0.0 s

8 16 24 32

0.2

0.4

ys

0.4 30

0.2

0.0 s

8 16 24 32

0.2

t

0.4

Fig. 5.8 Normal modes 30-32 for 32 masses on a string consisting of alternating unequal

masses, which are indicated by their size.

of motion by assuming separate solutions for the masses m1 and m2 , solutions with the

same periodicity but with dierent amplitudes. To explore this possibility, we relabel the

particles in pairs with the index s ranging from 1 to 16, and with vertical displacements s

and s corresponding to the sth masses m1 and m2 , respectively. Generalizing Eq. (5.2), the

hs+1 hN

x2 hs 1

h2 hs 1 xN 1

x1 xs xs+1 hN

h1 xs 1 xN

t

Fig. 5.9 Coordinates for alternating masses m1 and m2 on a massless string under tension T .

t

149 Sequences of different point masses on a string

d2 s T

m2 = ( s 2 s + s+1 ) , (5.79)

dt2 l

d2 s T

m1 2 = ( s 1 2 s + s ) (5.80)

dt l

where we recall that the beginning and end of the string are fixed so that 0 = 0 and

N+1 = 0 (for an even number of masses). To solve this system of coupled equations, we

begin by assuming trial solutions with the usual time dependence

s = (0)

s e

i!t

(5.81)

s = (0)

s e

i!t

. (5.82)

Substituting these trial solutions into Eqs. (5.79) and (5.80), we obtain

!

T (0) 2T T (0)

s + ! m2 (0)

2

s =0 (5.83)

l l l s+1

!

T (0) 2T T (0)

s 1 + !2 m1 (0)s =0. (5.84)

l l l s

We can solve this coupled system of equations following the same approach we used in

5.1.1 to solve the case where all the masses are the same. This time, however, we must

write down trial solutions for two displacement amplitudes (0) (0)

s and s . The numerical

solutions obtained in the previous section suggest spatially oscillating solutions where, for

a given normal mode, each set of particles m1 and m2 oscillate with the same wavelength.

Adopting complex notation once again, we assume trial solutions of the form

(0)

s = e

i(s )

(5.85)

1

(0)

s = e

i[(s 2) ]

= e i/2 i(s

e )

, (5.86)

where and are overall amplitude factors for the modes for masses m1 and m2 . Following

the same procedure we used in 5.1.1, we have once again explicitly included a phase factor

and thus take the amplitudes (0) and (0) to be real. The factor of s 12 in Eq. (5.86) takes

account of the fact that the m1 masses are displaced one half integer unit to the left of the

m2 masses. Applying the boundary conditions that 0(0) = (0) N+1 = 0 to these equations gives

0(0) = e i

=0 (5.87)

i[(N+ 12 )

(0)

N+1 = e

]

=0. (5.88)

Recalling that we are interested only in the real part of these equations, we can satisfy Eq.

(5.87) by setting = /2. Thus, Eq. (5.88) can be rewritten as

sin[(N + 12 )] = 0 , (5.89)

which is satisfied only if (N + 12 ) = n where n = 1, 2, 3, . . . . Therefore, the boundary

conditions once again restrict the values of :

n

n = , where n = 1, 2, 3, . . . , N , (5.90)

N + 12

t

150 Waves of oscillating particles

where we write n to emphasize that depends on the mode number n. Substituting these

results back into Eqs. (5.85) and (5.86), the amplitudes become

0 1

(0)

BBB ns CCC

ns = n sin @ B CA (5.91)

N + 12

0 1

(0)

BBB n(s 12 ) CCC

ns = n sin @ B CA . (5.92)

N + 12

where we have added the subscript n to explicitly designate the normal mode. Equations

(5.91) and (5.92) are essentially the eigenvectors that we obtained in the numerical solution

to this problem in the previous section. For example, the eigenvector for mode n would be

(0) (0) (0) (0) (0) (0)

aTn = (n1 , n1 , n2 , n2 , ..., nN , nN ). Notice, however, that these eigenvectors have not

been normalized nor have we determined the relationship between the overall amplitudes

n and n . We shall return to these issues after we determine the frequencies of the normal

modes.

Having satisfied the boundary conditions, we now need to see if our trial solutions satisfy

the equations of motion. Therefore, we substitute Eqs. (5.85) and (5.86) into Eqs. (5.83)

and (5.84), which gives

1 T h 1 i

m1 !2 ei[(s 2) ]

= ei((s 1) )

2 ei[(s 2) ]

+ ei(s )

(5.93)

l

T h 1 1 i

m2 !2 ei(s )

= ei[(s 2) ]

2 ei(s )

+ ei[(s+ 2 ) ]

. (5.94)

l

Canceling out the common exponential factors gives

T h i

m1 !2 = 2 + ei/2 + e i/2

(5.95)

l

2 T h i/2 i

m2 ! = e + e i/2 2 . (5.96)

l

Collecting together and rearranging terms we obtain

!2 !21 + !21 cos = 0 (5.97)

2

!2 cos + ! !22 = 0 .

2 2

(5.98)

2

where we have defined the characteristic frequencies !21 = 2T/m1 l and !22 = 2T/m2 l. This

system of equations has non-trivial solutions for the amplitudes and only when the

determinant of the coefficients of and is zero, that is,

=0 (5.99)

!22 cos /2 !2 !22

where we have used the identity sin2 (/2) = 1 cos2 (/2). The solution is readily obtained

t

151 Sequences of different point masses on a string

(s 1)

q

1200 12 + 22

p

1 = 2T /m1 l

1000

800

600 p

2 = 2T /m2 l

400

200

t

0 n

0 5 10 15 20 25 30

Fig. 5.10 Frequencies of normal modes vs. mode number for alternating sequence of

evenly-spaced masses m1 and m2 on a string. Analytical solution: open squares.

Numerical solution: solid circles.

2 s 3

1 666 4!21 !22 sin2 (n /2) 7777

!2n = !2 + !22 666641 1 777 . (5.101)

2 1 (!2 + !2 )2

1 2

5

The solution has two branches, !a and !o , corresponding to subtracting or adding the

discriminant in Eq. (5.101). The lower branch is called the acoustical" branch and the

upper is call the optical" branch, a nomenclature that comes from the analysis of the

vibrations of atoms in solids, an analysis essentially the same as we have followed here.

The two branches of the solution are plotted in Fig. 5.10. Notice that because we indexed

the particles in pairs, the mode number n runs from 1 to 16 for the analytical solution; for

each mode number n, there are two normal modes, one with frequency !a and the other

with frequency !o . Below the band gap, the frequencies obtained by numerically solving

the eigenvalue problem correspond to !a . Above the band gap, the frequencies obtained

numerically correspond to !o , but the frequencies are indexed dierently for the numerical

and analytical solutions: mode 1 in the optical branch of the analytical solution corresponds

to mode 32 of the numerical solution; mode 2 in the optical branch of the analytical solution

corresponds to mode 31 of the numerical solution, etc.

Exercise 5.2.1 In the limit that N 1, show that the highest frequency !max

a of the of

the acoustical branch and the lowest frequency !min

o of the of the optical branch are

t

152 Waves of oscillating particles

given by

r

2T

!max

a ' !2 = (5.102)

m2 l

r

2T

!min

o ' !1 = , (5.103)

m1 l

assuming that m1 < m2 . Note that there are the same frequencies we estimated pre-

viously for the modes just below and above the band gap by a physical argument

(see Eqs. (5.77) and (5.78)). Show that the maximum frequency !max

o of the optical

branch is given by

s !

q

max 2 2 2T 1 1

!o ' !1 + !2 = + . (5.104)

l m1 m2

Exercise 5.2.2 Starting from Eq. (5.101), show that

s !

2T 2n

!n ' . (5.105)

(m1 + m2 )l 2N + 1

for the acoustical branch when n N, which corresponds to the small kn or long

wavelength limit. For n N, Eq. (5.105) shows that !n increases linearly with n,

and thus linearly with kn as well.

Pulling together the various parts of our solution to the normal mode problem of alter-

nating masses on a string we have

N

X

ns (t) = n sin(n s) (An cos !n t + Bn sin !n t) (5.106)

i=1

N

X

1

ns (t) = n sin[n (s 2 )] (An cos !n t + Bn sin !n t) , (5.107)

i=1

where n is given by Eq. (5.90) and !n is given by Eq. (5.101). Note that we have expanded

the complex exponential term exp( i!t) so that the expressions for ns (t) and ns (t) are

entirely real. The amplitude ratio n /n for each mode is fixed by Eqs. (5.97) and (5.98),

while the overall amplitudes n and n of the various modes are determined by the initial

conditions.

Equations (5.106) and (5.107) can be written in terms the particle coordinates x s along

the string, just as we did at the end of 5.1.1. Here, however, the relevant distance is 2l, as

the pattern of alternating masses repeats itself every two particles rather than after every

particle as in 5.1.1. Thus we define the wave vector kn such that n s = kn x s , where x s =

2ls. The allowed values are

n

kn = , where n = 1, 2, 3, . . . , N . (5.108)

(2N + 1)l

The smallest value k1 corresponds to a disturbance with a wavelength 2/k1 that is twice

the length of the string, as shown in the top plots in Figs. (5.7) and (5.8). For the acoustical

branch, the masses m1 and m2 oscillate in phase, while for the optical branch, the masses

t

153 Sequences of different point masses on a string

m1 and m2 oscillate out of phase. The largest value of kN is approximately /2l, which

corresponds to a wavelength of 2/kn = 4l, is shown in the bottom two plots of Fig. 5.7.

masses on a string. Here we explore the normal modes of a random sequence of masses

on a string. Oddly enough, the normal modes that emerge are not completely random but

exhibit certain generic features, which have important consequences for the propagation

of waves in a wide variety of systems. Our simple system provides a convenient point of

entry into the fascinating physics of such systems.

First lets define what we mean by a random sequence of masses. By random, we mean

that we select without bias each mass along the string from a well-defined distribution of

masses. This distribution is defined at the outset of the selection process. For example, we

might choose masses from a uniform distribution of masses between some minimum mass

m1 and maximum mass m2 . Then the distribution function is:

8

>

>

> 1

>

< if m1 < m < m2

P(m) = > m 2 m1 (5.109)

>

>

>

:0 otherwise .

The probability that a particular mass is between m and m + dm is P(m) dm. Integrating

over all m should give unity

Z 1 Z m2

1

P(m)dm = dm = 1 , (5.110)

0 m1 m2 m1

as indeed it does.

Alternatively, we might choose masses from a Gaussian distribution of width about a

mean of m0 :

(m m0 )2 /2 2

P(m) = P0 e , (5.111)

where the constant P0 is fixed by the normalization condition that the integral over all m

must yield unity probability, as in Eq. (5.110). In either case you can imagine a process in

which masses are randomly selected out of a hat that contains a large number of particles

whose masses are distributed according to some well-defined distribution such as those

given above.

Here we consider a particularly simple distribution of masses: only two dierent masses,

m1 and m2 , but arranged in a random sequence. Problems at the end of the chapter explore

the normal modes of systems with other kinds of random distributions. For the present

case, we imagine that the probabilities of choosing a mass m1 or m2 are equal.

First we need to specify a random sequence for the two masses. We could do so by

flipping a coin: heads, choose m1 , tails, choose m2 ; repeat until you have a sequence of

32 masses. This is much too tedious, however, so we let a computer do it for us using an

algorithm known as a random number generator. We are not going to delve into how one

t

154 Waves of oscillating particles

0.4

1

ys

0.0

8 16 24 32

0.4 particle number s

2

0.0

ys

8 16 24 32

0.4 particle number s

0.4 16

ys

0.0

8 16 24 32

0.4 particle number s

0.4 17

ys

0.0

t

8 16 24 32

0.4 particle number s

Fig. 5.11 Normal modes for a random sequence of light particles (small circles) and heavy

particles (large circles). There are 32 total particles. Mode number is indicated in

upper right corner of each plot.

develops such an algorithm; you can read about that elsewhere.2 We simply assume that

we can find a random number generatora computer functionsuitable for our purposes.

Having generated a random sequences of masses, we need to solve the normal mode

problem. Once again, we turn to the numerical methods and the computer programs we

have developed. To solve the present problem, we need only to populate the diagonal ele-

ments of the mass matrix m with our random sequence of masses and then run the computer

program that diagonalizes the matrices and calculates the normal modes. We present the

results generated by such a numerical solution below.

Figure 5.11 shows the first two normal modes as well as normal modes 16 and 17 for

a random sequence of masses m1 and m2 where the total number of masses is 32. Here

we use the same physical parameters used in previous sections: m1 = 2.5 g, m2 = 10 g,

l = 4 cm, and T = 60 N. The lowest frequency modes look essentially the same as those

for the cases of a identical masses and alternating masses, respectively, as can be checked

by comparing Fig. 5.11 with Figs. (5.5) and (5.7). In the limit that the number of masses

N 1, the oscillation frequencies of the low n modes for all three cases is given by

r

T l n

!n = , for n N , (5.112)

m L

where m is the average mass of a bead on the string and L is the length of the string.

Exercise 5.2.3 Starting from Eqs. (5.38) and (5.105), show that Eq. (5.112) holds for

2 Numerical Recipes by Press et al. provides a useful introduction to random number generators.

t

155 Sequences of different point masses on a string

the cases of identical masses and alternating masses on a string. Equation (5.112)

also holds the in the case of a random

p sequence of masses. Why? Note that Eq.

(5.112) can be rewritten as !n = T/ (n/L) where = m/l is the mass per unit

length of the beaded string.

Modes 16 and 17 for our random sequences of masses are also shown in Fig. 5.11. In

contrast to the low-n low-freqency modes, these higher frequency modes do not resemble

the normal modes observed previously for the case of identical or for the case of alternating

masses. Moreover, modes 16 and 17 will look dierent from those depicted here for other

realizations of the system, that is, for other random sequences of the same masses m1 and

m2 . For example, Fig. 5.12 shows normal modes 16 and 17 for a dierent random sequence

of masses. This illustrates an important point: the particular shapes of the larger n normal

modes depends on the particular sequence of masses. By contrast, the low-n low-frequency

modes all look very similar and have similar normal frequencies.

Although the higher-n modes change significantly for dierent random sequences of

masses, they do exhibit some common features. In particular, a relatively small fraction

of the particles in modes 16 and 17 have large displacements, as shown in Figs. (5.11)

and (5.12). The region of large particle displacements is confined to a range of about 10-14

particles for modes 15 and 16. Such modes are said to be localized and are a generic feature

of all but the lowest frequency modes in systems with a random distribution of masses.

In general, as the mode number and frequency increases, the normal modes in a random

system become more and more localized. Figure 5.13 shows the three highest-n modes for

the original random sequence of masses used in Fig. 5.11 and we see that they are even

more localized than the n = 15 and 16 modes. We can characterize the spatial extent

the localizationof a normal mode quantitatively by averaging over the square of the

amplitudes of the particles in a given mode. Because the energy is proportional to the

square of the particle amplitude, we are weighting according to where energy of the mode

0.4 16

0.0

ys

8 16 24 32

0.4 particle number s

0.4

17

0.0

t

ys

8 16 24 32

0.4 particle number s

Fig. 5.12 Normal modes for a random sequence of light particles (small circles) and heavy

particles (large circles) that is different from that shown in Fig. 5.11. There are 32 total

particles. The mode number is indicated in the upper right corner of each plot.

t

156 Waves of oscillating particles

N

X

s = s y2s . (5.113)

s=1

P

Here we use normalized eigenvectors so that s y2s = 1 so no additional normalization is

needed. The spatial extent of a normal mode can then be characterized by

0 N 11/2

BBBX 2 2C

C

` = l B@ (s s) y s CCCA

B . (5.114)

s=1

The length ` gives the spatial extent of a normal mode, and is commonly known as the

localization length. In general, ` depends on the frequency of the mode and, to a lesser

extent, on the particular realization of the random sequence of masses. From the depiction

of normal modes in Fig. 5.11, we might expect that ` L for the low-n modes. By contrast,

for the From Figs. (5.12) and (5.13), we would expect ` L for the high-n high frequency

modes. Indeed, this trend is supported by Fig. 5.14(a), which shows how ` varies with

frequency for the system of 32 masses depicted in Figs. (5.11) and (5.13). As expected, `

decreases with ! at high frequencies but is relatively flat at low frequencies, with a value

of about one third of the systems sizeabout 10 particles out of 32. You can get a sense of

32

0.4

0.0

ys

8 16 24 32

0.4 particle number s

31

0.4

0.0

ys

8 16 24 32

0.4 particle number s

30

0.4

0.0

ys

8 16 24 32

t

0.4 particle number s

Fig. 5.13 High frequency normal modes for a the same random sequence of light particles and

heavy particles shown in Fig. 5.11. The mode number is indicated in the upper right

corner of each plot.

t

157 Sequences of different point masses on a string

how much ` can vary from one random sequence to another from the smaller gray points

that show ` vs. ! for 20 dierent realizations of a random sequence of 32 masses. Clearly,

the fluctuations in ` from one realization to another can be significant. The large variation in

` from one realization to another is due in part to the relatively small number of particles in

the system, 32 in this case. Figure 5.14(b) shows data for a string with 2048 masses, where

the fluctuations in ` are smaller and the trend for ` to decrease with increasing frequency

is quite clear. In fact, the results for this larger system are good enough to suggest that

` ! 2 for the high frequency modes, as shown by the gray line in Fig. 5.14(b). The

localization length ` ceases to increase as ` approaches the system size L, just as it did for

the smaller system.

Now that we know something about the normal modes of our random system, lets ex-

amine the normal frequencies. Figure 5.15 shows the normal frequencies as a function of

mode number for the system of 32 masses (dark circles). Interestingly, Fig. 5.15 reveals

that there appear to be a number band gaps, although the gaps are smaller than the single

gap we found for the case of alternating masses. Figure 5.15 also shows the normal fre-

quencies for a dierent random sequence of masses (light circles). The normal frequencies

of this second sequence of are very nearly the same for the low-frequency modes but show

significant dierences at the higher frequencies. The two realizations of the system have

band gaps at similar, but not exactly identical frequencies. This is not surprising as the two

sequences are dierent.

Exercise 5.2.4 Why are the normal modes of the two dierent realizations of a random

sequence of masses very nearly the same at low frequencies but fairly dierent at

other frequencies, particularly those in the mid-range? [Hint: See Exercise 5.2.3]

Examining Fig. 5.15 reveals that for both realizations of the system, there are no normal

modes in the frequency ranges from about 920 s 1 to 1130 s 1 and from about 1200 s 1 to

101

102

l

101

100

100

t

10 1 10 1

102 103 100 101 102 103

Fig. 5.14 Localization length vs. frequency. (a) N = 32. Red data shows ` for the sequence of

32 masses used in Fig. 5.11; black data are for 20 different realizations of sequences

of 32 masses. (b) Data for N = 2048 masses. Gray line shows ! 2 for reference.

t

158 Waves of oscillating particles

1350 s 1 . Is this a generic feature of all realizations of this system or just a coincidence for

these two realizations? To address this question, we need to examine the normal modes for

a large number of random systems that are statistically equivalent. By statistically equiva-

lent, we mean that for each realization, the sequence of masses is chosen randomly from

the same probability distribution. Such a collection of systems is called an ensemble, and

averaging some physical quantity of interest over this ensemble is called an ensemble av-

erage.

To answer the question about band gaps posed above, we calculate the normal frequen-

cies for a large number of realizations of this system and make a histogram of the fraction

of the total number normal frequencies between ! and !+ !, where ! is fixed to be some

convenient small frequency interval. We call frequency-dependent height of the histogram

the density of normal modes or, more commonly, the density of states D(!). Figure 5.16

shows the density of states calculated for an ensemble of 32 particles and for an ensemble

of 2048 particles. The density of states is essentially the same for the two systems, with

D(!) being more finely grained due to the larger number of modes. In both cases, there is

a remnant of the band gap just below !1 , the lowest frequency mode of the optical band

for a system of alternating spheres. Moreover, there is a surplus of states between !1 and

!t , the range of frequencies occupied by the acoustical branch for the case of alternating

masses. There are also states above !t , the maximum frequency normal mode for the case

of alternating spheres. Problem ... at the end of the chapter explores the physical origins of

these features of the density of states.

1600

1400

1200

1000

(rad/s)

800

600

400

200

t

0

8 16 24 32

mode n

Fig. 5.15 Normal modes for a random sequence of 32 spheres: dark circles, random sequence

of masses used in Fig. 5.11; light circles, random sequence of masses used in Fig.

5.12.

t

159 Sequences of different point masses on a string

masses is a quasiperiodic sequence of masses. A simple example of a 1-dimensional qua-

sicrystal can be constructed from a Fibonacci sequence. You may be familiar with a Fi-

bonacci sequence of integers: 1, 2, 3, 5, 8, 13, 21, 34, 55, 89, . . .. The sequence starts with

0, 1 with the rule that the next element fi of the sequence is generated from the sum of the

previous two elements: fi = fi 1 + fi 2 .

We can use the same rule to create a binary sequence of two letters, say A and B. Starting

with the sequences f1 =A and f2 =AB, f3 = f2 + f1 =ABA, f4 = f3 + f2 =ABAAB, f5 =

f4 + f3 =ABAABABA, etc., where in this case the +" operation denotes the concatenation

of two sequences. Among the many very interesting properties of the Fibonacci sequence

is that the pattern of letters A and B never repeats itself, somewhat like the digits of an

irrational number. However, the sequence of letters is not exactly random either, as we

shall see.

Here we consider a Fibonacci sequence of the same two masses m1 and m2 that we used

for the alternating sequence and for the random sequences of masses. This time, however,

we construct the sequence of masses the way we constructed the sequence of letters A and

(a) 2 1 t m

( )

(b) 2 1 t m

( )

t

0 200 400 600 800 1000 1200 1400 1600

Fig. 5.16 Density of states D(!) for a random sequence of masses m1 and m2 . (a) N = 32; (b)

N = 2048. Dark vertical lines indicate the frequencies !1 , !2 , and !t corresponding to

the frequencies of the top of the acoustic branch and the bottom and top of the optical

branch, respectively, for an alternating sequence of masses m1 and m2 . The frequency

!t corresponds to the maximum normal frequency consisting of all light masses m1 .

t

160 Waves of oscillating particles

1400

1200

1000

800

(rad/s)

600

400

200

t

0

8 16 24 32

mode n

Fig. 5.17 Frequencies of normal modes vs. mode number for a Fibonacci sequence of two

masses.

B above: the letters A and B correspond to the masses m1 and m2 , respectively. Instead

of constructing our systems of 32 or 2048 particles, the number of particles will be a

Fibonacci number fn . This means that the number of masses m1 and m2 will be fn 1 and

fn 2 , respectively. Lets start with N = 34, which means that there are 21 particles of mass

m1 and 13 particles of mass m2 arranged in the following sequence:

ABAABABAABAABABAABABAABAABABAABAAB. (5.115)

We solve the eigenvalue problem once again simply by substituting the sequence of

masses, in this case the Fibonacci sequence, into the mass matrix. Figure 5.17 shows the

resulting normal mode frequencies as a function of the mode number. Three obvious band

gaps are visible, just above modes 13, 21, and 26. Note that the number of modes between

each set of gaps in a Fibonacci number: 13, 8, and 5. Alternatively, we note that the three

gaps occur at 8, 13, and 21 modes, all Fibonacci numbers, counting right to left from the

highest frequency mode.

Figure 5.18 shows the two lowest frequency modes, which as in the cases for other

sequences of masses, are modes with wavelengths of 2L and L, where L is the length of

the string. Figure 5.18 also shows modes 13 and 14, the modes just below and just above

the lowest-frequency gap, respectively. In mode 13, the heavy masses m2 have the largest

amplitudes, while in mode 14, the light masses have the largest amplitudes with the large

particles remaining close to their equilibrium positions. It is this sudden change in the mass

that oscillates without a corresponding change in the wavelength that leads to the band gap.

As the normal mode frequencies increase, the oscillations become more localized, al-

though in a way that is quite distinct from the localized modes of the random system. As it

is difficult to appreciate the character of these high frequency modes for a system of only

t

161 Sequences of different point masses on a string

34 particles, we illustrate the high frequency modes for a system of 89 particles. Figure

5.19 shows the normal modes for the three highest frequencies of a two-mass system of

89 particles arranged in a Fibonacci sequence. Mode 87 is similar to localized modes in a

random system in that nearly all the particles with large displacements are near each other.

Modes 88 and 89 are more typical, however, of high frequency modes in a Fibonacci se-

quence. In both cases there are two sets of well-separated localized oscillations. Thus, one

cannot say that those modes are localized in the same sense that high frequency modes for

a random distribution of masses are localized. Unlike the system of identical particles or

the system of alternating masses, however, the oscillations do not appear to extend across

the entire system. Thus, the character of the normal modes of the Fibonacci sequences,

with their well-defined band gaps but quasi-localized high frequency modes, seems to be

somewhat intermediate between an ordered system of alternating masses and a completely

disordered system of randomly distributed masses.

The previous sections illustrate how the arrangement of the masses aects the normal

modes of the system. A regular alternating sequence of two masses produces a large band

gap as well as sinusoidal normal modes that span the entire system. By contrast, random

and quasiperiodic sequences produce smaller band gaps and, at least at the higher frequen-

cies, produce modes that are localized to regions that do not span the entire system. In this

section we explore a bit more the connection between how the masses are arrangedthe

structure of the systemand the properties of the normal modes.

0.4

1

ys

0.0

8 16 24 32

0.4 particle number s

2

0.0

ys

8 16 24 32

0.4 particle number s

0.4

13

0.0

ys

8 16 24 32

0.4 particle number s

0.4

14

0.0

t

ys

8 16 24 32

0.4 particle number s

Fig. 5.18 Normal modes for a Fibonacci sequence of two masses. The mode number is shown

in the upper right corner of each plot.

t

162 Waves of oscillating particles

0.4

89

0.0

ys

16 32 48 64 80

0.4 particle number s

0.4

88

0.0

ys

16 32 48 64 80

particle number s

0.4

0.4

87

t

0.0

ys

16 32 48 64 80

0.4 particle number s

Fig. 5.19 Normal modes for an 89-mass Fibonacci sequence of two masses. The mode number

is shown in the upper right corner of each plot.

A common way of characterizing the structure of a system is through the Fourier trans-

form, usually the power spectrum, of the density. One can choose the number density, mass

density, or whatever seems relevant to the system of concern. Here we choose to look at a

normalized power spectrum S (k) of the mass density, which we define as

1

S (k) = hk k i , (5.116)

N

where k is the discrete Fourier transform (see Eq. (7.99)) of the mass density

N

X

i(2/N)ks

k = (m s /m 1) e (5.117)

s=1

and m is the average mass of a particle in the system. We have written m s /m 1 instead of

m s /m to suppress the very large peak at k = 0 that arrises due to the fact that the average

mass along the chain is finite (i.e. not zero). The angle brackets in Eq. (5.116) designate an

ensemble average, which is relevant only in cases where there can be dierent statistical

realizations of the system, as for the random sequences of two masses we studied in 5.2.2.

The function S (k) is called the structure factor and is a measure of the extent to which a

system contains repeating spatial patterns in its structure.

Up to this point, we have been dealing with normal modes, which are standing waves.

These waves do not propagate from one place to another; each normal mode simply oscil-

lates in place in a pattern determined by the eigenvector for that particular mode. In this

section, we explore traveling waves, that is, waves that move from one place to another.

t

163 Traveling waves

S(k)

80

40

100 50 0 50 100

S(k)

40

20

100 50 0 50 100

S(k) 2

t

100 50 0 50 100

k

Fig. 5.20 Structure factor S (k) for (top) alternating masses, (middle) a Fibonacci sequence of

masses, and (bottom) a random sequence of masses. The black line is the ensemble

average of 64 realizations. The gray line shows S (k) for a single realization.

Traveling waves are of immense practical importance as they carry energy with them and

thus can move or transport energy from once place to another. Ocean waves, for example,

transport mechanical energy over thousands of kilometers. An earthquake under the sea

floor can generate a wave visible on the ocean surface that can deliver enormous energy to

distance places, often with disastrous consequences. On a less dramatic, but no less impor-

tant scale, sound waves can deliver acoustic energy created by an orator or an orchestra to

the farthest reaches of a vast auditorium or concert hall.

Our purpose in this section is to explore the propagation of waves, in particular wave

pulses. We are concerned with the velocities at which waves and wave pulses propagate

and how the width of a wave pulse can change as it propagates.

When a wave runs into a barrier, it is generally reflected, sometimes refracted, and often

partially absorbed. While reflection, refraction, and absorption are fascinating phenomena,

they represent complications that, for the moment, we wish to avoid. Therefore, we shall

defer these phenomena so that we can concentrate on wave propagation without these

complications.

To avoid the problem of reflection, we will work with systems without any borders,

systems that have no beginning or end. For one dimensional systems this is simple to

realize: instead of working with a linear system confined between two walls, we connect

t

164 Waves of oscillating particles

the two ends of our one dimensional system to form a ring: voil, a system without borders,

a system with no beginning and no end.

Once again, we consider masses on a string, as in the previous sections, but this time with

the ends connected in a ring. Now you may wonder how we can maintain tension in such a

string. One way would be to spin the string around like a wheel to create a centrifugal force.

Or alternatively, we could make the string" out of a wire of some convenient stiness, but

still with a mass small compared to the masses of the beads we distribute along it so that

we could ignore the mass o the wire. However you imagine accomplishing this feat, we

will consider the string or wire to be characterized by a tension T with masses distributed

equal distances l around its circumference. In the end, considering wave propagation on a

ring is merely as useful artifice for avoiding reflections and other end eects.

We begin with the simplest case, N identical masses m on a string with the masses

indexed from s = 1 to N just as we did in 5.1. Once again, the transverse displacement of

the sth mass is y s . As every bead is free to move, we no longer have the boundary conditions

y0 = yN+1 = 0 that ensure that the ends of the string are stationary. However, because the N

beads are arranged on a ring, y s = y s+N , for the sth bead and the (s + N)th bead are the same

bead. The requirement that y s = y s+N is referred to as a periodic boundary condition. The

equations of motion are essentially the same as those developed in 5.1 with the equation

of motion for the sth bead given by Eq. (5.2)

T

my s = (y s 1 2y s + y s+1 ) . (5.118)

l

Because the masses are on a ring, the first mass with s = 1 is next to the last mass with

s = N. Thus for s = 1 and s = N, Eq. (5.118) becomes

T

my1 = (yN 2y1 + y2 ) (5.119)

l

T

myN = (yN 1 2yN + y1 ) . (5.120)

l

Once again, we can recast this system of equations as a matrix equation of the form

m y = k y , (5.121)

where m and y are the same as those given below Eq. (5.7). The periodic boundary condi-

tions embodied in Eqs. (5.119) and (5.120) lead to a small change in the stiness matrix k ,

which remains a tri-diagonal matrix except that two new entries appear at the upper right

and lower left corners of the matrix:

0 1

BBB 2 1 1CCC

BBBB 1 2 1 0 CCCC

BBB CCC

BBB . . CCC

BBB 1 2 . CCC

T BBBB B CCC

k = BB . .

.. .. .. . CCC .

l BBBB CCC

CCC

BBBB .. CCC

BBB . 2 1 CCC

BBB 0 1 2 1CCCCC

BB@ A

1 1 2

t

165 Pulse propagation and dispersion

Evolution of a pulse

The solutions given by Eqs. (5.106) and (5.107) are sums of the normal modes of the prob-

lem, which are standing waves that oscillate between the two fixed endpoints or boundaries

of the string.

Thus far we have investigated the normal modesthe wavesof discrete masses uni-

formly distributed along a massless string under tension. Here we begin to explore the

normal modes of a simple string with a mass per unit length of . Building on what we

have learned thus far, a straightforward way to proceed is to consider the limit where the

number of discrete masses N goes to infinity while the mass per unit length remains

constant. The length L = Nl of the string remains constant, which means that the distance

l between masses becomes negligibly small in the limit N ! 1.

We start by considering the equation of motion for a typical particle of mass m along a

string, which is given by Eq. (5.2) and rewritten here for convenience:

T

my s = (y s 1 2y s + y s+1 ) . (5.122)

l

In the limit that l ! 0, the quantity in parentheses on the right hand side of Eq. (5.122) is

proportional to @2 y/@x2 . To see how this comes about, consider three closely spaced points

along an arbitrary continuous function y(x), as depicted in Fig. 5.21. The first derivative at

x s can be expressed as

dy ys ys 1 y s+1 y s

= lim = lim , (5.123)

dx x!0 x x!0 x

where x s remains fixed while x s 1 and x s+1 move as x ! 0. The two expressions for

the derivative in Eq. (5.123) are called the backward dierence and forward dierence,

respectively, and approximate the derivative either a half step backwards or forwards from

x s for a finite step width x. In the limit x ! 0, they yield the same result.

Using the backward and forward dierence expressions as approximations for the deriva-

tives at a half step forwards and a half step backwards, we can obtain an expression for the

second derivative at x s as follows

y s+1 y s ys ys

d2 y y0s+1/2 y0s 1/2 x x

1

dx2 x!0 x x!0 x

y s+1 2y s + y s 1

= lim , (5.125)

x!0 ( x)2

t

166 Waves of oscillating particles

ys+1

ys

ys 1

t

x

xs 1 xs xs+1

which is proportional to right hand side of Eq. (5.122). Dividing Eq. (5.122) by m and then

setting = m/l and l = x gives

T (y s 1 2y s + y s+1 )

y s = lim . (5.126)

l!0 ( x)2

We recognize the limit on the right hand side as the partial second derivative of y with

respect to x. In the limit that l ! 0, the displacement y becomes a function of x as well as

t so both derivatives in Eq. (5.127) are partial derivatives. Thus, Eq. (5.126) can be written

as a second order partial dierential equation

@2 y T @2 y

= (5.127)

@t2 @x2

The ratio T/ has units of velocity squared and thus Eq. (5.127) is usually written as

@2 y @2 y

2

= 32 2 , (5.128)

@t @x

p

where 3 = T/. Equation (5.128) applies quite generally to strings. However, it also de-

scribes a number of other kinds of waves, including electromagnetic waves (light), as well

as sound, provided the transverse displacement y is replaced by the appropriate quantity

such as electric field or molecular displacements. Because of its generality, Eq. (5.128) is

usually referred to as the wave equation.

Solutions to the wave equation can be written in a very general form

where f is any continuous function that can be dierentiated twice. For such a function

@2 f (x 3t) 2 00 @2 f (x 3t)

= 3 f (x 3t) , = f 00 (x 3t) , (5.130)

@x2 @x2

which clearly satisfies the wave equation. The function f (x 3t), whatever the shape of

f (x), is simply a waveform that moves at constant velocity v in the +x direction. Similarly,

f (x + 3t) moves at constant velocity v in the x direction.

t

167 The continuum approximation

t

Fig. 5.22 Gaussian wave packet propagating to the right in a nondispersive medium. The black

curve shows the wave at t = 0 and the red curve shows the same wave a time t later.

For simple sinusoidal waves, the function f (x, t) is just a sine or cosine wave, which for

a sine wave we typically write as

!

= A sin k x t = A sin [k(x 3t)] , (5.132)

k

where v = !/k. Having written A sin(kx !t as A sin [k(x 3t)], we see that it is simply a

sinusoidal waveform moving to the right at a speed of 3 = !/k.

Next we consider the propagation of a more interesting wave, a Gaussian wave packet,

which we first encountered in 7.2.1. To obtain a traveling wave packet, we simply let

x ! x vt in Eq. (7.42). This gives

(x 3t)2 /2a2

f (x, t) = f (x 3t) = B e cos k0 (x 3t) . (5.133)

The fact that f (x, t) = f (x 3t) guarantees that Eq. (5.133) satisfies the wave equation.

Of course, you can verify this by explicitly performing the time and spatial derivatives of

f (x, t). Figure 5.22 shows the f (x, t) at t = 0 and at a slightly later time, which illustrates

how the wave packet moves along the +x axis. Notice how the wave packet moves to the

right undistorted as it propagates.

5.5.1 Dispersion

For many waves encountered in nature the velocity of the wave depends on frequency

(and thus on wavelength). For example, red light propagates through most glass at a faster

velocity than does blue light. Glass prisms exploit this fact to separate white light into the

various colors of the visible spectrum. Another example comes from the problem we have

spent most of this chapter investigating, the propagation of waves on a string with evenly

spaced masses. Recall from Eq. (5.37) that the frequency ! is related to the wave vector k

by

1

! = 2!0 sin kl , (5.134)

2

t

168 Waves of oscillating particles

where we treat the frequency and wave vector as continuous variable and thus have dropped

the subscript n. Because ! does not depend linearly on k, the velocity, which is given by

3 = !/k, is not a constant but depends on k (or !).

The relationship between ! and k is called the dispersion relation. Because 3 = !/k, the

velocity 3 will not be constant but instead will depend on k (or !) whenever the dispersion

relation is a nonlinear function of k. This has important consequences for the propagation

of pulses, which we explore below.

The fact that 3 is not a constant means that solutions of the form y(x, t) = f (x 3t)

are not longer valid, as we assumed that v was a constant in obtaining those solutions. Of

course, for any single wavenumber k, the velocity is a constant (since k is not changing),

so sinusoidal waveforms are still valid solutions of the wave equation. This suggests that

we analyze the wave equation in terms of its Fourier components. We start by rewriting the

wave equation as

@2 y 2

2@ y

3 =0, (5.135)

@t2 @x2

Next we write y(x, t) in terms of its spatial Fourier transform

Z 1

1

y(x, t) = Y(k, t) eikx dk . (5.136)

2 1

That is, at each value of t, we take the spatial Fourier transform of y(x, t)

Z 1

Y(k, t) = y(x, t) e ikx dx . (5.137)

1

Writing y(x, t) in terms of its Fourier transform Eq. (5.136), its derivatives become

Z 1

@2 y 1

= k2 Y(k, t) eikx dk (5.138)

@x2 2 1

Z 1 2

@2 y 1 d Y(k, t) ikx

= e dk . (5.139)

@t2 2 1 dt2

Notice that d2 Y(k, t)/dt2 is a full derivative since Y is not a function of x. Substituting these

expressions into Eq. (5.135) and collecting the terms under one common integral gives

Z 1 2 !

1 d Y(k, t) 2 2

+ k 3 Y(k, t) eikx dk = 0 . (5.140)

2 1 dt2

The only way the Fourier integral can be zero for all values of t is by setting the quantity

inside the parentheses is zero. Doing so leads to the equation

d2 Y(k, t)

+ !2 (k) Y(k, t) = 0 , (5.141)

dt2

where we have used ! = k3, emphasizing that in general the frequency is a function of k by

writing ! = !(k). Equation (5.141) represents a significant simplification for the analysis

of the wave equation, as it has reduced a second order partial dierential equation in two

variables, x and t, to an ordinary second order dierential equation in one variable, x. Thus

we need only solve the ordinary dierential equation for Y(k, t), from which we can then

obtain the full solution for y(x, t) using Eq. (5.136)

t

169 The continuum approximation

In fact, Eq. (5.141) is quite easy to solve as it is just the usual equation for a simple

harmonic oscillator. Thus, its solution is given by

where Y(k, 0) is determined by the initial conditions. Substituting Eq. (5.142) into Eq.

(5.136), we arrive at the full general solution to the wave equation

Z 1

1

y(x, t) = Y(k, 0) ei[kx!(k) t] dk . (5.143)

2 1

Equation (5.143) can be interpreted very simply. Given a waveform at time t = 0, the

waveform at later times is obtained by summing (integrating) the Fourier components of

the wave at t = 0 and propagating the sum forward in time using each Fourier components

frequency as given by the dispersion relation !(k).

Lets illustrate how Eq. (5.143) works with a couple of examples. We start with a very

simple example where the wave is made up of just two sinusoidal waveforms with dierent

wave vectors and frequencies:

y2 (x, t) = A sin(k2 x !2 t) . (5.145)

These two waves each have a single Fourier component, so Eq. (5.143) tells us that we

can sum these two Fourier components and simply propagate them forward in time each at

their own frequency. Summing these two waveforms gives

!t) (5.146)

where we have used sine addition and subtraction formulas and where

k = k1 k2 , ! = !1 !2 (5.147)

1 1

k = 2 (k1 + k2 ) , =

! 2 (!1 + !2 ) . (5.148)

t

Fig. 5.23 Superposition of two sinusoidal waves of similar wavelength and frequency but with

different velocities. The resultant waveform is characterized by two velocities: (1) The

phase velocity 3 p , which is the speed at which the blue waveform moves, and (2) The

group velocity 3g , which is the speed at which the envelope shown as a gray dashed

line moves. For dispersive waves, these two velocities are not the same: 3 p , 3g .

t

170 Waves of oscillating particles

Figure 5.23 shows a plot of Eq. (5.146) at t = 0. The full waveform, shown in blue, consists

of a sinusoidal wave, sin(kx !t) modulated by a sinusoidal envelope, cos( kx !t),

shown as dashed gray lines. The velocity of the sharp peaks and crests of the blue waveform

sin(kx !t)

is called the phase velocity and is given by 3 p = !/ k. By contrast, the velocity

of the (gray) envelope cos( kx !t) is called the group velocity and is given by 3g =

!/ k. For a nondispersive wave where ! = k3, 3 p = 3g . However, for a dispersive wave,

where !(k) is nonlinear in k, 3 p , 3g , so the peaks and crests for the blue waveform travel

at a dierent velocity than the envelope.

Next we consider a sinusoidal waveform with a Gaussian envelope, a waveform we

considered previously in 7.2.2, as it illustrates nearly all the generic features of pulse

propagation for dispersive waves. We start by specifying the waveform at t = 0:

x2 /2a2

y(x, t = 0) = B e cos k0 x . (5.149)

We already calculated the Fourier transform of this wave packet in 7.2.2. From Eq. (7.55)

we have

1 1 2 2 1 2 2

Y(k, 0) = B a e 2 (k k0 ) a + e 2 (k+k0 ) a (5.150)

2

According to Eq. (5.143), we need to specify the dispersion relation !(k) along with Y(k, 0)

in order to fully determine the time development of y(x, t). A simple model that shows the

salient features of dispersion is provided by the dispersion relation

!

k2

!(k) = !0 1 + 2 , (5.151)

2

so that the integral in Eq. (5.143) becomes

Z 1

1 a 1 2 2 1 2 2

!0 (1+ 12 k2 /2 )t]

f (x, t) = e 2 (k k0 ) a + e 2 (k+k0 ) a ei[kx dk . (5.152)

2 1 2

The integrals in Eq. (5.152) can be performed exactly by completing the squares in the

exponents. The result is the real part of (**redo integral**)

p 8 2 ! t 2 2 9

(t) >

>

< 12 (t) k0 x (kk00/) 0

i[k0 x !(k0 )t]

1

2 (t)

k0 x+(k0 /)2 !0 t >

>

i[k0 x+!(k0 )t] =

f (x, t) = > e a

e + e k0 a

e > ,

2 > : >

;

(5.153)

where (t) = [1 + i!0 t/(a)2 ] 1 . The result corresponds to two identical Gaussian wave

packets, one traveling in the negative direction and the other traveling in the positive direc-

tion, as shown in Fig. 5.24.

The first exponential in Eq. (5.153) is the Gaussian envelope of the wave packet. Its

center position is given by the condition k0 x (k0 /)2 !0 t = 0 or x = (k0 /2 )!0 t. Dier-

entiating, we obtain dx/dt = (k0 /2 )!0 , the group velocity 3g of the wave packet. In the

previous example consisting of just two frequency components we found that the group

velocity was given by 3g = !/ k. Here, where the wave packet consists of a continuous

distribution of frequencies, the expression for the group velocity becomes

d!

3g = . (5.154)

dk k=k0

t

171 The continuum approximation

x

50 0 50 100 150 200 250 300 350 400

t

Fig. 5.24 Propagation of Gaussian wave packet with dispersion. The blue trace shows wave

packet at t = 0 where it consists of right and left propagating parts. Red traces show

wave pack at later times, showing how the wave packet gets wider with time. The gray

bars have a length of 2a(t).

3g = k0 !0 /2 , (5.155)

Figure 5.24 shows that as the wave packet propagates to the right, its width increases.

We can obtain a quantitative expression for the width of the wave packet from Eq. (5.153),

from which it is evident that a2 / (t) is the square of the width of the Gaussian envelope.

Because (t) is complex, we take its square modulus, from which we obtain the following

expression for the time dependent width of the wave packet

1/2 " ! t 2 #1/2

a2 0

a(t) = =a 1+ 2 2 . (5.156)

(t) a

The spreading of the pulse arises because there is a spread of wave vectors k in the pulse,

which leads to a spread in the group velocities. Thus some parts of the wave packet travel

faster than others, leading to an increase the the width of the packet. Note that in the

previous example where we were adding together only two sinusoidal waves, there was

only a single group velocity 3g = !/ k precisely because there were only two sinusoidal

waves. In that case, therefore, there is no spreading of the wave packet as it propagates.

It is evident not only that the Guassian wave packet spreads as it propagates but also

that its amplitude decreases. This is as it must be. The energy associated with the wave

packet becomes more spread out in space so the amplitude must decrease as energy is

proportional to the square of the amplitude. In this example there is no dissipation, that is,

no loss of energy as the wave propagates. If there were, the amplitude of the wave packet

would decrease even more rapidly.

This brings us to another important point: energy is generally transmitted by wave pack-

ets at the group velocity. Obviously, as the wave packet spreads, some of the energy prop-

agates a bit faster while other propagates a bit slower. Nevertheless, energy transmission

is generally associated with the group velocity. This has some interesting consequences

for electromagnetic wave propagating in dispersive media where the velocity is a strong

t

172 Waves of oscillating particles

x

80 100 120 140 160 180 200 220

t

Fig. 5.25 Gaussian wave packet as its center passes x = 150 for three times t1 < t2 < t3 , where

t1 is in red, t2 in solid blue, and t3 in dashed blue. The envelope is shown in gray

dashed lines.

function of frequency. In some cases, it is possible for the phase velocity to exceed the

speed of light c while the group velocity remains below the speed of light. Because energy

(and information) is not propagated faster than the speed of light, the principle of special

relativity is not violated. Nevertheless, it is interesting that situations exist where 3 p > c

while 3g < c.

The phase velocity of the Gaussian wave packet is given by

0 1

!(k0 ) !0 BBB k02 CCC

3p = = @B1 + 2 AC , (5.157)

k0 k0 2

which means that the ratio of the phase to the group velocities is

3 p 1 2

= + 2 . (5.158)

3g 2 k0

For the propagating pulse shown in Fig. 5.24, /k = 2 so 3 p /3g = 4.5. This means that the

crests and valleys in the wave packet move 4.5 times faster than the Gaussian envelope.

The dierence in the phase and group velocities is illustrated in Fig. 5.25 where we show

the wave packet at three times in rapid succession as the center of the wave passes x =

150. From the figure it is clear that the crests and valleys of the wave packet are moving

considerably faster than the Gaussian envelope. During the period of time that the crests

and valleys advance an appreciable fraction of a wavelength, the envelope scarcely moves

at all.

While we have considered a particular example of dispersion in this section, the phe-

nomena illustrated here are generic. When the phase velocity is not constant but depends

on frequency (or wavelength), the group velocity is dierent from the phase velocity and

in general a wave packet spreads as it propagates. This has important consequences for

all kinds of waves, and none more important than for quantum mechanical waves. For ex-

ample, as we shall see, the dispersion relation for an electron (and any other particle with

mass) is highly nonlinear with the frequency increasing quadratically with wave vector.

Thus, electrons and other particles tend to spread out in space as they propagate.

t

173 Problems

Problems

5.1 Consider a system of five particles of mass m evenly spaced on a massless string

fixed at both ends and stretched in a straight line to a tension T . The middle mass

is then pulled up so that the system has the initial configuration shown in the figure

below. The particles are initially all at rest.

0.06

0.04

0.02

0.00

0 1 2 3 4 5 6

0.02

(a) At time t = 0, the masses are released and begin to oscillate. Starting from Eq.

(5.31), find the subsequent motion of all five the particles. Hint: This problem

boils down to finding the ten coefficients {Bn } and {Cn } in Eq. (5.31). You can

set all but three of the ten coefficients to zero by considering the initial con-

ditions, including the symmetry of the initial configuration with respect to the

symmetries of the five normal modes indexed by n.

(b) Using your solution from part (a), plot the configurations of the particles at sev-

eral dierent times as the system oscillates. Does the system ever return to its

initial configuration? If so, when? If not, why?

5.2 Show that the limit that 1, Eq. (5.101) yields the following two solutions

r

!1 !2 3 T

!a = + O( ) = + O(3 ) (5.159)

2(!21 + !22 )1/2 2(m1 + m2 )l

s !

2 2 1/2 2 2T 1 1

!o = (!1 + !2 ) O( ) = + O(2 ) . (5.160)

l m1 m2

Therefore, as discussed in the text, the acoustic branch has a linear dispersion rela-

tion in the limit of small or long wavelengths. By contrast, the dispersion relation

for the optical branch is nearly independent of frequency and thus highly dispersive.

5.3 Starting from Eqs. (5.97) and (5.98) derive the following expressions for the normal

mode amplitude ratios

!22 !2a

n /n = , acoustical branch (5.161)

!22 cos(/2)

!21 !2o

n /n = , optical branch , (5.162)

!21 cos(/2)

t

174 Waves of oscillating particles

and plot the results using the expressions for !a and !o given by Eq. (5.101). Set

N = 32 for your plots and explain briefly in words what the plots signify for n = 1

and at the band gap.

5.4 Show by writing y(x, t) in terms of its temporal Fourier transform Y(x, !) (dierent

from Y(k, t)) that you can obtain a dierent equation

d2 Y(x, !)

+ k2 (!) Y(x, !) = 0 , (5.163)

dx2

where in this case Y(x, !) is the temporal Fourier transform of y(x, t).

5.5 Consider a linear string of masses all with the same mass m connected to each other

by springs whose spring constants alternate between k1 and k2 as shown in the fig-

ure below. The masses are constrained to move only in the horizontal direction (no

transverse motion).

(a) Write down the system of equations of motion for the masses and then express

them in the standard matrix form m x = k x. What are the mass and stiness

matrices m and k ? Note: You will need to write down separate equations of

motion for for masses with even and odd values of s.

(b) Setting m = 10 g, k1 = 3 N/m, and k2 = 6 N/m, solve for the eigenvectors and

eigenvalues of the system for the case of N = 16 masses using the numerical

routines available on the NYUClasses course web site. Plot x s vs. s for modes

1-4, 7-10, and 13-16 (i.e. the eigenvectors) Plot all the normal frequencies as a

function of mode number from lowest to highest mode number. You should find

a band gap between modes 8 and 9. Compare your results to the case of two

dierent masses but identical springs discussed in the text. Explain physically

why there is a band gap by considering the modes just below and just above the

gap.

6 Strings

6.3 Interfaces

175

t

176 Strings

Problems

6.1

7 Fourier Analysis

The waves we encounter in nature, whether they are the ripples on the surface of a pool

of water or the sounds of a jazz ensemble, are seldom characterized by the simple single-

frequency and single-wavelength sine and cosine functions introduced in elementary ac-

counts of waves. The sound waves produced by even the simplest musical instrument have

a much more complex harmonic content, meaning they are made up of waves simultane-

ously oscillating a number of dierent frequencies. It is this rich harmonic content that

makes an A-note played on a violin sound dierent from the same" A-note played on a

trumpet. Take a look at the waveforms shown in Fig. ??, which were taken from a violin and

a trumpet. The two waveforms have the same overall period, but look rather dierent, as

you might expect given the very dierent quality of the sounds the two instruments make.

The waveforms dier in their harmonic content, that is, in the spectrum of frequencies that

make up the sound we hear when each instrument is played. Fourier analysis attempts to

characterize mathematically, and rather precisely, just how the waveforms and harmonic

content of the trumpet and the violin dier.

The ripples on the surface of a pool of water present another kind of waveform, one that

is clearly a wave that oscillates up and down, but one that has a finite extent: it starts at

some some position in space, grows bigger, reaches a maximum, and then decays away.

This waveform is more-or-less sinusoidal, except its amplitude grows and decays in space.

If instead of taking a snapshot of the spatial undulations of the waveform, you monitored a

single point on the surface of the water, you would see it start to oscillate up and down in

time as the wave passed by. Once again, the amplitude of the waveform initially grows, this

time in time, reaches a maximum, and then decays away. How do we describe the harmonic

content this waveform? The waveform superficially resembles the beats we encountered

in our study of two coupled oscillators in Chapter 4. There we found that two coupled

oscillators could produce beats if both modes were simultaneously excited and if the two

modes had very similar frequencies. This provides a clue as to the harmonic content of

the waveform of the water surface wave: perhaps it consists of two or more sinusoidal

waves of nearly the same frequencies. This is in fact the case. The precise mathematical

description of such waves is once again the subject of Fourier analysisin this case Fourier

transforms.

In this chapter, we study Fourier series and Fourier transforms. Fourier series can de-

scribe virtually any periodic function. However, Fourier series can only describe periodic

functions of infinite extent, that is, functions that repeat forever either in space or time. To

describe wave pulses, that is, waveforms of finite extent in space or time, we need a more

general approach, which involves Fourier transforms. In 7.1 we introduce Fourier series

and in 7.2 we introduce Fourier transforms.

177

t

178 Fourier Analysis

Computational physics, which involves encoding information in digital form, puts fur-

ther demands on our mathematical description of waves. To address these demands, we

develop discrete Fourier transforms (DFTs) in 7.3 as well as algorithms for calculating

DFTs, including the famous fast Fourier transform (FFT) algorithms.

The basic idea of Fourier series is that virtually any periodic function can be represented

as an infinite series of sine and cosine functions. Fourier series are useful not only as

mathematical tools but also as a means for understanding the harmonic makeup of complex

non-sinusoidal periodic waveforms.

In physics, we generally encounter functions that are either periodic in space or in time or

in both space and time. In this section, we explicitly deal with functions that are periodic in

space. However, we could equally equally well work in the temporal domain. Ultimately,

we will study functions that are periodic in both space and time.

Figure 7.1 shows two examples of periodic functions. Mathematically, we can express

the fact that a function f (x) is periodic in space by writing f (x + L) = f (x). We call L

the periodicity (or wavelength) of the function. Such a function can quite generally be

represented by an infinite series of sine and cosine waves:

X 1 X 1

1

f (x) = a0 + an cos kn x + bn sin kn x (7.1)

2 n=1 n=1

where kn = 2n/L and n = 1, 2, 3, ... . For n = 1, k1 = 2/L and the sine and cosine waves

have a wavelength of L. For n = 2, 3, ..., the sine and cosine waves have successively shorter

wavelengths L/n corresponding to the harmonics of the n = 1 sine and cosine waveforms,

as shown in Fig. 7.2.

For a given periodic function f (x) there is a unique set of coefficients an and bn that,

when used in Eq. (7.1), yield a Fourier series that converges to the desired periodic func-

tion. We can obtain the coefficients using a simple procedure. First, we multiply both sides

(a) 6 (b) 2

x x

t

8 6 4 2 2 4 6 8 4 3 2 1 1 2 3 4

6 2

Fig. 7.1 Periodic wavefunctions: (a) sawtooth wave with a period L = 4, (b) triangle wave with

a period L = 2.

t

179 Fourier series

n=1 1 n=1 1

1 1

n=2 1 n=2 1

1 1

n=3 1 n=3 1

t

L/2 L/2 L/2 L/2

1 1

Fig. 7.2 Cosine (blue) and sine (red) waves for the first three terms n = 1 to 3 in a Fourier

series. The solid lines indicate the interval L/2 x L/2.

of Eq. (7.1) by cos km x, where km = 2m/L and m can be any positive integer, and then

integrate over the interval [ L/2, L/2]:

Z L/2 Z L/2

1

f (x) cos km x dx = a0 cos km x dx

L/2 2 L/2

1

X Z L/2

+ an cos kn x cos km x dx

n=1 L/2

X1 Z L/2

+ bn sin kn x cos km x dx . (7.2)

n=1 L/2

The first integral on the right hand side of Eq. (7.2) is zero since cos km x has positive and

negative parts in equal measure on the interval ( L/2, L/2). The third integral on the right

hand side of Eq. (7.2) is also zero since the product sin kn x cos km x is an odd function in-

tegrated over an interval ( L/2, L/2) that is symmetric about x = 0. Finally, the second

integral is zero unless km = kn . If km = kn , then the integral is just 1/2, the average of

cos2 kn x, times L, the length of the interval of integration. Mathematically, we can summa-

rize the integrals over the sine and cosine functions by

Z L/2

L

cos kn x cos km x dx = nm (7.3)

L/2 2

Z L/2

sin kn x cos km x dx = 0 (7.4)

L/2

Therefore, the only term in the sum over an on the right hand side of Eq. (7.2) that is

t

180 Fourier Analysis

Z L/2 1

X Z L/2

f (x) cos km x dx = an cos kn x cos km x dx

L/2 n=1 L/2

X1

L

= an nm

n=1

2

L

= am . (7.5)

2

We can solve this to obtain as expression for am . Once we have obtained the result, however,

the index label is arbitrary, so we can write the result using the index n:

Z L/2

2

an = f (x) cos kn x dx . (7.6)

L L/2

An equation for bn can be found using a similar approach by multiplying both sides of Eq.

(7.1) by sin km x and integrating over the interval [ L/2, L/2]. The result is

Z L/2

2

bn = f (x) sin kn x dx , (7.7)

L L/2

Finally, an expression for a0 can be obtained simply by integrating both sides of Eq. (7.1)

over the interval [ L/2, L/2]. The sine and cosine terms integrate to zero giving:

Z L/2

2

a0 = f (x)dx . (7.8)

L L/2

Before proceeding, we point out that we could have performed the integrals in Eqs. 7.6

7.8 over any interval of length L and still have obtained the same results for the coefficients

an and bn . For example, we could have used the interval [0, L]. Try it out! Because any

waveform described by a Fourier series is periodic with a periodicity of L, the choice of

interval is up to you.

We are now in a position to use Eq. (7.1), together with Eqs. 7.67.8 for the Fourier

coefficients an and bn , to determine the Fourier series for a few example waveforms.

Lets start by finding the Fourier series for the waveforms shown in Fig. 7.1. Finding the

Fourier series for a periodic waveform f (x) is, according to Eq. (7.1), simply a matter of

determining the coefficients an and bn , which we can do using Eqs. 7.67.8. To calculate the

integrals in Eq. (7.1), we need to specify the function f (x) only in the interval of integration

[ L/2, L/2], i.e. one period. For the waveform in Fig. 1(a), the interval of integration is

t

181 Fourier series

(a) 6 (b) 6

x x

8 6 4 2 2 4 6 8 8 6 4 2 2 4 6 8

6 6

(c) 6 (d) 6

x x

t

8 6 4 2 2 4 6 8 8 6 4 2 2 4 6 8

6 6

Fig. 7.3 Fourier series (black) for sawtooth function (red) including: (a) n = 1 term, (b) n 2

terms, (c) n 3 terms, and (d) n 4 terms.

[ 2, 2] ( L = 4) where the function is f (x) = 3x. The coefficients are then given by

Z Z

2 L/2 2 2

a0 = f (x) dx = 3x dx = 0 (7.9)

L L/2 4 2

Z !

2 L/2 2nx

an = f (x) cos dx (7.10)

L L/2 L

Z nx

1 2

= 3x cos dx = 0 (7.11)

2 2 2

Z L/2 !

2 2nx

bn = f (x) sin dx (7.12)

L L/2 L

Z nx

1 2

= 3x sin dx (7.13)

2 2 2

12 12

= 2 2 [n cos(n) sin(n)] = ( 1)n+1 , (7.14)

n n

where we used the fact that cos(n) = ( 1)n and sin(n) = 0 for n = 1, 2, 3, ... to evaluate

Eqs. (7.11) and (7.14). Substituting the values for an and bn back into Eq. (7.1), we obtain

the Fourier series for f (x):

X1

12 nx

f (x) = ( 1)n+1 sin (7.15)

n=1

n 2

!

12 x 1 1 3x

= sin sin x + sin ... . (7.16)

2 2 3 2

Figure 7.3 shows f (x) together with a series of plots that include successively more terms in

the Fourier series for f (x). As more terms are included in the sum, the series approximates

the periodic function better and better. In fact, including only those terms up through n = 4

results in a surprisingly good approximation of the sawtooth waveform.

The sawtooth waveform shown in Fig. 7.3 is an odd function. That is f (x) = f ( x).

Therefore, the Fourier series for f (x) must be odd, term by term. This means that no even

term, that is, no term proportional to cos kn x can contribute to the sum since cosine is an

t

182 Fourier Analysis

(a) 6 (b)

x 0 x

8 6 4 2 2 4 6 8 1 2 3 4 5 6 7 8

t

6 6

Fig. 7.4 Fourier components that make up the Fourier series plotted in Fig. 7.3: (a) sine waves

with correct amplitudes (and phases) (b) amplitudes bn = ( 1)n+1 12/n of sine waves.

even function and even one such term would spoil the odd symmetry of f (x). Therefore

we can set the coefficients an of the even cosine terms in Eq. (7.1) to zero. By noting

this symmetry, we can avoid having to explicitly perform the integrals in Eq. (7.11) to

determine the coefficients an . We only need to determine the coefficients bn of odd terms,

sin kn x, as only they can contribute to the Fourier sum.

Figure 7.4(a) shows the sine waveforms that were added together to make up the saw-

tooth wave. Note that they consist of the fundamental periodicity L and all its harmonics

L/n. Figure 7.4(b) shows the amplitudes bn of the sine waves that make up the sawtooth

wave. Note that because the power in a wave is proportional to the amplitude squared, the

power in each frequency is proportional to b2n . More generally, the power in each frequency

is proportional to a2n + b2n , but in this (special) case it reduces to b2n because an = 0 for all

n. We develop the subject of the power in a wave described by a Fourier series in 7.1.4.

The sets of coefficients {an } and {bn } give a quantitative measure of the harmonic content

of a particular waveform. They tell how much of a waveform is made up of sine and cosine

waves of a particular wavenumber kn , or equivalently of a wavelength n = 2/kn . As

noted in the introduction to this chapter, it is the harmonic content, that is the particular

collection of wavenumbers kn and their amplitudes an and bn , that gives a particular musical

instrument its characteristic sound.

To reinforce these ideas we look at another example. Consider the triangle waveform in

Fig. 7.1(b):

(

4x + 2 , 1x0,

f (x) = (7.17)

4x + 2 , 0x1.

Note that it is an even function: f ( x) = f (x). Because it is an even function, the coeffi-

cients of the sin kn x terms are zero, i.e. bn = 0. Therefore, we only need to determine the

coefficients an , which are given by Eq. (7.6):

Z

2 L/2

an = f (x) cos kn x dx

L L/2

Z 0 ! Z 1 !

2nx 2nx

= (4x + 2) cos dx + ( 4x + 2) cos dx

1 2 0 2

(

16/(n2 2 ) , if n is odd;

= (7.18)

0 , if n is even;

t

183 Fourier series

1

X 16

f (x) = 2 2

cos(nx) (7.19)

n=1,3,5,...

n

" #

16 cos(3x) cos(5x)

= 2 cos(x) + + + ... . (7.20)

9 25

The Fourier series developed above to describe periodic waveforms in space can also be

used to describe periodic waveforms in time. Doing so is as simple as replacing the spatial

variable x with time t and the wavenumber kn with the angular frequency !n where

2n

!n =

, (7.21)

T

and T is the period of the waveform. With these changes, we write

X1 X 1

1

f (t) = a0 + an cos !n t + bn sin !n t . (7.22)

2 n=1 n=1

In this case the Fourier series is a sum over sines and cosines of dierent temporal fre-

quencies !n . Because of this correspondence between spatial and temporal Fourier series,

people often refer to the wavenumber kn that appears in the spatial Fourier series as a

spatial frequency, in analogy to the temporal frequency !n . Thus, spatial frequency and

wavenumber are simply two names for kn .

The sine and cosine functions that we use in Fourier series are often referred to as ba-

sis functions. The integrals over their products, like those given in Eqs. (7.3) and (7.4),

are known as scalar products and express the orthogonality the the sine and cosine basis

functions.

This language of basis functions and orthogonality is an extension of the concepts and

language we use when speaking of a Cartesian coordinate system. In a Cartesian coordinate

system, for example, we can express the coordinate of any point a as linear combination

of the three orthogonal basis vectors, e x , ey , and ez , the unit vectors along the x, y, and z

directions, respectively. Similarly, we can express any periodic function with period L as

linear combinations of the sin kn x and cos kn x basis functions (where kn = 2n/L). Two

Cartesian vectors a and b are orthogonal if their dot product is zero, that is if

3

X

ab= ai bi = 0 , (7.23)

i=1

readily adapt the definition of the dot (or scalar) product to 100 component vectors (or any

other number). Consider, for example the two 100 component vectors A = cos 2i/100 and

t

184 Fourier Analysis

B = sin 2i/100, where i runs from 1 to 100. Generalizing the dot product written above

to 100-dimensional vectors, we write the dot product of these two vectors as

100

X

AB= Ai Bi . (7.24)

i=1

We can increase the number of components to be as large as we wish and the scalar product

written above still makes sense. In the case of Fourier series, we have simply increased

the number of of components to infinity, because there is an infinite number of values of

x between L/2 and L/2. In so doing the sum over the various components in Eq. (7.24)

becomes an integral. For the Fourier series we have defined, there is also an infinite number

of basis functions cos kn x and sin kn x (where kn = 2n/L) because n runs from 1 to infinity.

With these considerations in mind, the scalar products of the sine and cosine basis functions

are:

Z L/2

L

cos kn x cos km x dx = nm (7.25)

L/2 2

Z L/2

L

sin kn x sin km x dx = nm (7.26)

L/2 2

Z L/2

sin kn x cos km x dx = 0 (7.27)

L/2

If the scalar product, thus defined, is zero, then we say that the sine and cosine basis vec-

tors are orthogonal. Equations (7.25)-(7.27) show that all the dierent basis Fourier basis

functions are indeed orthogonal, as their scalar products are always zero. Only when per-

forming a scalar product of a basis vector with itself does the scalar product yield a finite

answer, similar to the behavior of the Cartesian basis vectors e x , ey , and ez .

Rereading the first paragraph of this section, you should now understand the terminology

introduced there about basis functions, scalar products, and orthogonality.

Exercise 7.1.1 Write a computer program to compute the sum in Eq. (7.24) for A =

cos 2i/100 and B = sin 2i/100 where i runs from 1 to 100. Do you get zero?

Repeat the calculation for AA and BB. Do the numerical values of the dot products

you obtain make sense? Briefly explain.

The power in a waveform is proportional to the square of the wave. For a periodic function

f (x) represented by a Fourier series, the energy per period (wavelength) is proportional to

t

185 Fourier series

Z

2 L/2

P/ | f (x)|2 dx (7.28)

L L/2

Z X1 X1 2

2 L/2 1

= a0 + an cos kn x + bn sin kn x dx (7.29)

L L/2 2 n=1 n=1

0 Z 1 Z L/2 1 Z 1

2 BBBB 1 L/2 2 X X L/2 CC

= @B a dx + 2

an 2

cos kn x dx + 2

bn sin kn x dxCCCA

2

(7.30)

L 4 L/2 0 n=1 L/2 n=1 L/2

0 1 1

1

2 BBBB 1 2 1 X 2 1 X 2 CCCC

= B@ La0 + L a + L bC (7.31)

L 4 2 n=1 n 2 n=1 n A

1

a20 X

= + a2n + b2n . (7.32)

2 n=1

The step from Eq. (7.29) to Eq. (7.30) follows from the orthogonality of the sine and

cosine functions, which makes all the integrals on the interval [ L/2, L/2] not involving

the square of a sine or cosine function equal to zero. The integrals over cos2 kn x and sin2 kn x

are simply L/2, the average value of cos2 kn x and sin2 kn x, which is 1/2, times the length

of the integration interval, which is L. The equality between Eqs. (7.28) and (7.32),

Z 1

2 L/2 a2 X

| f (x)|2 dx = 0 + a2n + b2n (7.33)

L L/2 2 n=1

is known as Parsevals theorem. In this context it means that the power Pn in the nth Fourier

component is proportional to the sum of the square of the Fourier amplitudes. For mathe-

matical convenience we make the proportionality in Eq. (7.28) an equality and define the

power spectrum as

Pn = |an |2 + |bn |2 , (7.34)

which gives the power in the nth Fourier component for n 1. Because the n = 0 term

is not associated with any oscillations, we generally do not include it in the power spec-

trum. Thus, the power contained in the part of the frequency spectrum characterized by the

wavevector kn is given by Pn , aside from a numerical constant that depends on the type of

wave being describedsound, electrical, etc. The total power per period is given by sum-

ming Pn over all n. For example, the power in the sawtooth wave depicted in Fig. 7.1(a)

with a Fourier series given by Eq. (7.19) is

X1 !2 1

12 144 X 1

Ptot = = 2 = 24 . (7.35)

n=1

n n=1 n2

Performing the same calculation for the triangle wave depicted in Fig. 7.1(b) gives

X1 !2 1

16 256 X 1 8

Ptot = 2 2

= 4 4

= . (7.36)

n=1,3,5,...

n n=1,3,5,...

n 3

t

186 Fourier Analysis

Exercise 7.1.2 What are the wavelengths of the first three spatial frequencies kn with

non-zero power in the sawtooth and triangle waves depicted in Fig. 7.1? What frac-

tion of the total power in each waveform is contained in (i) the longest wavelength?

(ii) the longest two wavelengths with non-zero power? (iii) the longest three wave-

lengths with non-zero power?

Fourier series can also be expressed in exponential form by exploiting the identity ei =

cos + i sin . In this case, we write the Fourier series as

X1

f (x) = An eikn x , (7.37)

n= 1

where once again kn = 2n/L and L is the period of the waveform. In general, the coeffi-

cients An are complex. Complex coefficients are required to correctly capture the phase and

parity (oddness or evenness) of f (x). The coefficients can be determined using a procedure

similar to the one we used to obtain the coefficients of the sine and cosine Fourier series.

In this case, we multiply both sides of Eq. (7.37) by e ikm x and integrate over one period:

Z L/2 Z L/2 0B X 1

1

C

BBB ikn x C

f (x) e ikm x

dx = B@ An e CCCA e ikm x dx

L/2 L/2 n= 1

1

X Z L/2

= An ei(kn km )x

dx

n= 1 L/2

X1

ei(n m) e i(n m)

= An L

n= 1

2i(n m)

X1

sin[(n m)]

= An L

n= 1

(n m)

X1

= An L nm

n= 1

= Am L , (7.38)

where we have used the fact that sin[(n m)] = 0 for n , m (n and m are integers) and

the fact that lim!0 sin / = 1 (here = (n m)).1 Thus, Eq. (7.38) gives

Z

1 L/2

An = f (x) e ikn x dx (7.39)

L L/2

It is left as an exercise for the reader to show that the coefficients An can be expressed in

terms of those for the sine and cosine series

8 1

>

>

> (a|n| + ib|n| ) , n < 0 ;

< 21

>

An = >

> a0 , n=0; (7.40)

> 21

>

:

2 (an ib n ) , n > 0 .

1 You can show that lim!0 sin / = 1 by writing the Taylor series

sin / = ( 3 /3! + 5 /5! ...)/ = 1 2 /3! + 4 /5! ...), which for ! 0 is equal to 1.

t

187 Fourier transforms

Equations (7.37) and (7.39) constitute a completely equivalent alternative form to the sine

and cosine Fourier series expressed by Eq. (7.1) and Eqs. (7.6)-(7.8). The power spectrum

for waves described by a complex exponential Fourier series is given by

Pn = |An |2 = An An . (7.41)

The primary importance of the exponential form of Fourier series is that the exponential

form is readily generalized to describe pulsesthat is, functions of finite extent. This is

the subject of Fourier transforms, which are discussed in the next section.

When we encounter waves in nature, they have a finite spatial extent (and a finite temporal

duration). They are not infinitely periodic extending from x = 1 to x = 1, as is the case

for the waves described by the Fourier series encountered in 7.1. Nevertheless, it is still

useful to have some means of describing the harmonic content of such waves. As we shall

see, the mere fact that real waves have a finite extent or a finite duration means that more

than one frequency is required to describe their harmonic content.

Gaussian envelope:

x2 /2a2

f (x) = B e cos k0 x . (7.42)

Figure 7.5(a) shows a plot of Eq. (7.42) for the case that there are many oscillations in the

pulse, i.e. a k0 1 . The periodicity (wavelength) is approximately 2/k0 = 2/3 ' 2.1.

But unlike a cosine wave, this wave pulse has a finite duration of about 5a = 40. We cannot

use a Fourier series to evaluate the harmonic content of such a wave pulse because it does

not repeat. We could, however, use a Fourier series to evaluate the harmonic content of the

periodic wave in Fig. 7.5(b). It has been made periodic by repeating the wave pulse in Fig.

7.5(a) at intervals of x = 100. The exponential form of the Fourier series for this wave is

given by Eq. (7.37) with the Fourier coefficients given by Eq. (7.39) with L = 100 and f (x)

given by Eq. (7.42).

If our goal is to describe the harmonic content of a pulse such as Eq. (7.42), the value of

L we choose is arbitrary so long as we choose L to be much greater than the extent of the

wave pulse, i.e. so long as L a. As we make L larger and larger, our approximation of

f (x) should become better and better. In the general case, we let L go to infinity so that we

can describe a wave pulse of arbitrary (but finite) duration. To this end, we rewrite the sum

t

188 Fourier Analysis

(a) f (x)

x

20 10 10 20

(b) f (x)

t

200 100 100 200

Fig. 7.5 (a) Cosine wave with Gaussian envelope: k0 = 3 and a = 8 (arrow in background

extends from a to a). (b) The same cosine wave with Gaussian envelope as in (a)

but repeated at intervals of x = 100.

in Eq. (7.37) as

1

X

f (x) = lim An eikn x (7.43)

L!1

n= 1

X1 L !

2

= lim An ei 2nx/L n

L!1

n= 1

2 L

X1

LAn i 2nx/L 2 n

= lim e

L!1

n= 1

2 L

Z 1

1

= F(k) eikx dk , (7.44)

2 1

where limL!1 2 n/L ! dk (note that n = 1 so the sum is not changed by its intro-

duction) and we have defined the function F(k) = limL!1 LAn . Just as the amplitude An

determines how strongly waves of spatial frequency kn = 2n/L contribute the the Fourier

sum of Eq. (7.43), so the function F(k) determines how strongly waves of spatial frequency

k contribute the the Fourier integral of Eq. (7.44). Using the expression for An given by Eq.

(7.39), we obtain an expression for F(k):

Z !

1 L/2

F(k) = lim LAn = lim L f (x) e ikn x dx (7.45)

L!1 L!1 L L/2

Z 1

= f (x) e ikx dx . (7.46)

1

The function F(k) given by Eq. (7.46) is the Fourier transform of f (x). The inverse trans-

form of F(k) is given by Eq. (7.46). We rewrite this Fourier transform pair again for refer-

t

189 Fourier transforms

ence:

Z 1

dk ikx

f (x) = e F(k) (7.47)

1 2

Z 1

ikx

F(k) = dx e f (x) . (7.48)

1

For notational purposes, it is useful to define the Fourier transform operator and its inverse

Z 1

F {...} dx e ikx {...} (7.49)

1

Z 1

F 1 {...} dx eikx {...} . (7.50)

1

Thus, we can write F { f (x)}, denoting the Fourier transform of f (x), which is equal to F(k).

Similiarly, we can write the inverse Fourier transform as F 1 {F(k)}, which is equal to f (x).

We shall make use of this notation in 7.2.4.

Equation (7.47) is the Fourier transform of f (x), which is a function of the spatial vari-

able x. Often we want to take the Fourier transform of a function g(t) that is a function

of time t. Mathematically, it is no problem to simply replace the spatial variable x with

the temporal variable t in Eqs. (7.47) and (7.48). The Fourier transform variable k then

becomes the angular frequency !. Everything works fine if we make the simple substitu-

tion of t for x and of ! for k. However, in physics it is conventional to write the temporal

Fourier transform pairs with the opposite sign convention to that used in Eqs. (7.47) and

(7.48) above. That is, we write

Z 1

d!

g(t) = G(!) e i!t (7.51)

1 2

Z 1

G(!) = dt g(t) ei!t . (7.52)

1

This choice of sign convention presents no difficulty mathematically, but it might seem odd

to use dierent sign conventions for spatial and temporal variables. The reason for doing

so to do with the fact that a wave propagating in the positive x direction is given by ei(kx !t) ,

that is, with a positive sign in front of the kx term and a minus sign in front of the !t term.

This is the origin of the sign convention for spatial and temporal Fourier transforms often

used in physics, and we adopt it here.

We note that the Fourier transform pairs given by Eqs. (7.47) and (7.48) and by Eqs.

(7.51) and (7.52) do not incorporate the factor of 2 symmetrically. On the other hand, if

we rewrite the temporal Fourier transform pairs, in terms of the conventional frequency

= !/2, Eqs. (7.51) and (7.52) become

Z 1

g(t) = d e i 2t G() (7.53)

1

Z 1

G() = dt ei 2t g(t) . (7.54)

1

Written in terms of the frequency instead of the angular frequency !, the temporal Fourier

transform pairs appear as completely symmetric transforms of each other. Remembering

this fact can help you remember which transform includes the factor of 2.

t

190 Fourier Analysis

F(k)

t

k

4 3 2 1 0 1 2 3 4

Fig. 7.6 Fourier transform of a Gaussian wave packet. The Gaussian wave packet is given by

Eq. (7.42) and is plotted in Fig. 7.5(a). Its Fourier transform, plotted above, is given by

Eq. (7.55).

Having introduced the formalism of Fourier transforms, lets use it to find the Fourier

transforms of a few functions. We start with Eq. (7.42), a Gaussian modulated cosine wave.

We use Eq. (7.48) to calculate the Fourier transform of Eq. (7.42):

Z 1

F(k) = f (x) e ikx dx

Z 1

1

2 2

= B e x /2a cos k0 x e ikx dx

1

r

1 (k k0 )2 a2 1 2 2

= Ba e 2 + e 2 (k+k0 ) a (7.55)

2

The result, plotted in Fig. 7.6, is just what you might have expected: there are peaks in

the Fourier transform at k = k0 corresponding to the primary frequency of the cosine

wave. You may wonder why there are two peaks, one centered about k0 and the other

centered about k0 . The reason is readily understood by recalling the complex exponential

representation of cosine:

cos k0 x = 12 eik0 x + e ik0 x (7.56)

A straightforward calculation of the Fourier transforms of eik0 x and e ik0 x reveals the the

1 2 2 1 2 2

former is proportional to e 2 (k k0 ) a while the latter is proportional to e 2 (k+k0 ) a . Since

cos k0 x contains both eik0 x and e ik0 x , its Fourier transform has peaks centered at both k0

and k0 .

The two peaks are Gaussian functions centered about k = k0 having a frequency width

of approximately

p a 1 . To be more precise, the half width of the peaks at 1/e of the peak

height is 2/a. This means that the range of wavenumbers k required to form a pulse of

length x ' a is approximately the reciprocal of the pulse duration a 1 .

In the limit that k0 goes to zero, cos k0 t approaches unity so that the pulse shown Fig.

7.5 becomes a simple Gaussian without any oscillations. That is, in this limit

B x2 /2a2

f (x) = p e , (7.57)

2

t

191 Fourier transforms

1 2 2

F(k) = Ba e 2k a . (7.58)

Thus, we see that the Fourier transform of a Gaussian function with a width of x ' a is

another Gaussian function with a wavenumber width of k ' a 1 . That is, the characteristic

widths of a function and its Fourier transform are inversely proportional:

1

x/ . (7.59)

k

Please note that the proportionality constant in the above relation is not the same for all

Fourier transform pairs but depends on the functional form of f (x). We will have more to

say about this inverse proportionality in the next section, where we make the statement in

a more general and mathematically precise way. For now, suffice it to say that it is a very

general and an extremely important property of pulses and their Fourier transforms. It is

known as the uncertainty relation or uncertainty principle for waves.

Another common function is the exponential:

|x|/a

g1 (x) = e . (7.60)

With a little algebra you can show that its Fourier transform is given by

Z 1

2a

G1 (k) = e |x|/a e ikx dx = (7.61)

1 1 + (ka)2

|x|/a

g2 (x) = e cos k0 x , (7.62)

is

a a

G2 (k) = + (7.63)

1 + (k + k0 )2 a2 1 + (k k0 )2 a2

The functions G1 (k) and G2 (k) are Lorentzian functions centered about wavenumbers of

zero and k0 , respectively. Thus, the Fourier transform of an exponential is a Lorentzian.

Exercise 7.2.1 Perform the integral in Eq. (7.61) to calculate the Fourier transform of

g1 (x) = e |x|/a . Hint: Start by explaining why the following equality is true:

Z 1 "Z 1 #

e |x|/a e ikx dx = 2 Re e x/a e ikx dx (7.64)

1 0

(

1 , |x| a

g3 (x) = (7.65)

0 , |x| > a ,

t

192 Fourier Analysis

g3 (x)

(a) 1

0.5

x

3 2 1 0 1 2 3

(b) 0.8 G3 (k)

0.4

t

k

20 15 10 5 5 10 15 20

Fig. 7.7 (a) Square pulse and (b) its Fourier transform.

Z 1

ikx

G3 (k) = g3 (x) e dx

1

Z a

ikx

= e dx

a

ikx a

e eika e ika

2 sin ka

= = =

ik a ik k

= 2a sinc ka , (7.66)

where we have used the definition sinc x sin x/x. Equation (7.66) is plotted in Fig. 7.7(b).

Note that since sin x and x are both odd functions, their quotient sin x/x is even.

Unlike the Fourier transforms of the Gaussian and Lorentzian, which decay monotoni-

cally with increasing k, the Fourier transform of the square pulse oscillates about zero as

it decays away. The origin of the oscillationscalled ringing"is the abrupt change in

amplitude of the square pulse. Very high spatial frequencies k are required in the Fourier

transform F(k) in order to reproduce these sharp features. A useful rule of thumb is that

smooth features in a pulse produce smooth Fourier transforms while sharp features in a

pulse produce Fourier transforms with oscillations at large k.

In the previous section we pointed out that the characteristic width of a pulse and its Fourier

transform are inversely proportional to each other. This idea is summarized by Eq. (7.59).

We would like to make the statement about the relationship between the width x of a

pulse and its Fourier transform k more precise. To do so we need a precise definition of

width of a function.

t

193 Fourier transforms

| f (x)|2

P f (x) = R 1 . (7.67)

1

| f (x)|2 dx

As the power in the waveform f (x) is proportional to | f (x)|2 , P f (x) dx gives the fractional

power in the wave between the positions x and x + dx. We then define the average position

of the wave as

Z 1

hxi = x P f (x) dx , (7.68)

1

where the angular brackets h...i to denote the average. Thus, we define the average position

weighted according to what fraction of the waves power is located near any given point.

Similary, we can define the mean square width or variance of x as

Z 1

Var(x) = (x hxi)2 |P f (x)|2 dx . (7.69)

1

p

x= Var(x) , (7.70)

or

Z 1

2 2

( x) = h(x hxi) i = (x hxi)2 |P f (x)|2 dx . (7.71)

1

The quantity x gives us a quantitative measure of how far in space a pulse extends from

its average position hxi. A similar equation applies to k, as is illustrated in the worked

examples below.

h(x hxi)2 i = hx2 i hxi2 (7.72)

Lets use Eqs. (7.71) and (7.72) to find x and k for some of the Fourier transform

pairs we considered above. We start with the Gaussian pulse defined by Eq. (7.57) and its

Fourier transform Eq. (7.58). We calculate ( x)2 using Eq. (7.72), noting first of all that

hxi = 0 because f (x) is centered about x = 0. Therefore, in this case

( x)2 = hx2 i hxi2 = hx2 i (7.73)

Z 1

= x2 P f (x) dx (7.74)

1

Z 1 2

1 B x2 /2a2

= R1 x2 p e dx (7.75)

1

| pB2 e x2 /2a2 |2 dx 1 2

= 12 a2 . (7.76)

Thus we find

a

x= p (7.77)

2

t

194 Fourier Analysis

We can calculate k following a similar procedure. First we see that hki = 0 because F(k)

is symmetric about x = 0. Therefore

Z 1

= k2 PF (k) dx (7.79)

1

Z 1

1

= R1 k2 |F(k)|2 dk (7.80)

2

|F(k)| dk 1

1

Z 1

1 1 2 2

= R1 1 2 2

k2 |B a e 2k a )|2 dk (7.81)

|B a e 2 k a )|2 dk 1

1

1

= 2 , (7.82)

2a

Thus we find

1

k= p (7.83)

2a

Combining the expressions for x and k we obtain

1

x k= (7.84)

2

Equation (7.84) is the precise statement of the uncertainty principle for a Gaussian wave

packet. Similar expressions can be obtained for other wave packets, for example, the expo-

nential pulse given by Eq. (7.60) and its Fourier transform Eq. (7.61). However, for wave

packets other than the Gaussian wave packet, the numerical value of the produce x k

is generally greater than 1/2. The Gaussian wave packet yields the minimum value of the

product x k. Therefore, the most general expression of the uncertainty principle is writ-

ten as an inequality to cover all cases:

1

x k . (7.85)

2

Exercise 7.2.3 Find x and k for the exponential pulse given by Eq. p (7.60) and its

Fourier transform Eq. (7.61). Show that the product of x k = 1/ 2. Is this consis-

tent with the uncertainty relation Eq. (7.85)?

The Fourier transforms of the functions given in Eqs. (7.42) and (7.57) illustrate a more

general property of Fourier transforms. If F(k) is the Fourier transform of f (x), i.e. if

F { f (x)} = F(k), then it is straightforward to show that

Equation (7.111) is known as the frequency shifting theorem. Rewriting cos k0 x in complex

exponential form, the frequency shifting theorem can be used to obtain Eq. (7.42) from Eq.

t

195 Discrete Fourier Transforms

(7.57) or to obtain Eq. (7.63) from Eq. (7.61). We can obtain a similar relation relating a

spatially shifted function to its Fourier transform

i kx0

F { f (x x0 )} = F(k)e . (7.87)

Below, we state without proof a few other useful theorems about Fourier transforms:

( n )

d f (x)

F = (ik)n F(k) (7.88)

dxn

dn F(k)

F {xn f (x)} = in (7.89)

dkn

F { f (x) g(x)} = F(k) G(k) , (7.90)

where f (x) g(x) denotes the convolution of f (x) with g(x)

Z 1

f (x) g(x) f (x) g(x ) d . (7.91)

1

P(k) = |F(k)|2 = F(k) F (k) . (7.92)

The power between spatial frequencies k and k + dk in the spectrum is equal to (aside from

some dimensional constant) P(k) dk. The total power is obtained by integrating P(k) dk

over all k.

Fourier transforms provide a very powerful means for understanding and analyzing wave-

forms and other fluctuating functions. When using a computer or dealing with experimen-

tal data, the functions encountered are necessarily represented by discrete data sets defined

over a finite range rather than continuous functions defined over an infinite range. To ex-

tend the tools of Fourier transforms to such data sets, it is useful to develop a discrete

Fourier transform or DFT. In this case, the function that is to analyzed is represented by a

finite set of evenly-spaced data points a distance 4x apart, as illustrated in Fig. 7.8.

To see how the DFT works, we start by recalling the expressions for the continuous

Fourier transform pair:

Z 1

F(k) = f (x) e ikx dx (7.93)

1

Z 1

1

f (x) = F(k) eikx dk . (7.94)

2 1

We wish to form the discrete equivalent of the continuous Fourier transform where the

function f (x) is uniformly sampled over a range of length L at discrete intervals x = L/N,

where N is taken to be an even integer. The length L is chosen to correspond to the range

of x where f (x) is non-zero; ideally f (x) = 0 outside this range. To capture all the spectral

t

196 Fourier Analysis

information in f (x), the interval x should be chosen so there are at least two points per

cycle of the highest spatial frequency present in the waveform.

In formulating the DFT, the convention is to choose the interval over which the function

f (x) is sampled to be 0 x L with

L

x=n x, n = 0, 1, 2, ..., N 1, x= . (7.95)

N

Note that the interval is not symmetric about x = 0, in contrast to the definition of the

Fourier transform. This choice means has some practical implications for representing

waveforms with DFTs that we return to later. For the moment, we simply note that this

choice is the standard way that DFTs are defined.

Because there are only N independent samples of f (x), there are only N independent

values of F(k). The smallest absolute value of k is k = 2/L so that the discrete values of

k are given by

2

k=l k, l = N/2, ..., N/2 , k= . (7.96)

L

The range of spatial frequencies required to discretely represent f (x) must still include

both negative and positive values. This follows from the definition of the Fourier transform,

which requires positive and negative spatial frequencies irrespective of the range of x over

which f (x) is non-zero.

f (x) (a)

f3

f2

f1 fN 2

f0 fN 1

x0 x1 x2 x3 x

xN 2 xN 1

(b)

n=N 1

(c)

t

n=1

Fig. 7.8 (a) Function f (x) sampled at evenly-spaced intervals for the discrete Fourier

transform (DFT). (b) Shortest wavelength (smallest spatial frequency) represented by

sampling at an interval x. (c) Longest finite wavelength represented by DFT defined

over an interval of length N x.

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197 Discrete Fourier Transforms

Before forging forward with the mathematical development, we also note that the spatial

frequencies that are included in the DFT range from k = k = 2/L to k = (2/L)(N/2) =

2/[L/(N/2)]. These two extreme values of wave vectors we are considering correspond

to (i) a wavelength equal to the entire interval L over which the data are defined (see Fig.

7.8(b)), and (ii) to a wavelength that is two data points long L/(N/2). This is so because

it takes a minimum of two data points to represent the oscillations of a sine wave, one

point for the negative trough and another point for the positive peak (see Fig. 7.8(c)). If

the function in question has contributions at higher frequencies, they will not be properly

captured by the DFT because the function is not sampled at sufficiently short intervals.

With these definitions, the discrete version of the product kx is given by

! !

2 L 2

kx = (l k) (n x) = l n = ln (7.97)

L N N

Defining fn = f (n x), the discrete version of the Fourier transform becomes

Z 1 N 1

X

F(k) = f (x) e ikx dx F(l k) = fn e i (2/N)ln x . (7.98)

1 n=0

The discrete Fourier transform (DFT) is defined as the last sum in Eq. (7.98) less x:

N 1

X

i(2/N)ln

Fl = fn e . (7.99)

n=0

With this definition of Fl , we can also write down the discrete version of the inverse Fourier

transform:

Z 1 N/2

X1

1 1

f (x) = F(k) eikx dk f (n x) = Fl ei (2/N)ln , (7.100)

2 1 l= N/2

N x

1 1 2 1 2 1

dk = = . (7.101)

2 2 L 2 N x N x

The inverse discrete Fourier transform (iDFT) is defined as the last sum in Eq. (7.101) less

1/ x:

N 1

1 X

fn = Fl ei(2/N)ln . (7.102)

N l=0

Notice that in going from Eq. (7.101) to Eq. (7.102) we have changed the sum over l from

N/2 to N/2 1 to 0 to N 1. There is no real problem with this change since Fl is periodic

in l with a period of N, as can be seen from Eq. (7.102). You just have to keep in mind that

the non-negative spatial frequencies appear in the first half (from l = 0 to l = N/2 1)

followed by the negative spatial frequencies (from l = N/2 to l = N 1) in the second half,

as detailed in Table 7.1. The non-intuitive ordering of Fourier components in the DFT is

done, in part, so that the two sums in the definitions for the DFT and the iDFT, Eqs. (7.99)

and (7.102), respectively, both run from 0 to N 1. Another reason for this choice is that

DFTs are useful primarily in computer applications where it is convenient for indices such

t

198 Fourier Analysis

Fourier transform.

l kl l kl

0 0 N/2 ( N/2) k

1 k N/2 + 1 ( N/2 + 1) k

2 2 k N/2 + 2 ( N/2 + 2) k

.. .. .. ..

. . . .

N/2 2 (N/2 2) k N 2 2 k

N/2 1 (N/2 1) k N 1 k

as n and l to run from 0 (or 1) to to some finite number, in this case N 1 (or N). It is no

use protesting that you might prefer to do it another way. This is the standard that everyone

uses.

Discrete Fourier transforms can be calculated very efficiently using a computer algo-

rithm called the Fast Fourier Transform or the FFT. There are several variants of the FFT

algorithm, the oldest of which dates back to to Gauss in 1805, although the basic algo-

rithm has been independently rediscovered by a number of people since that time. Many

mathematical computer libraries implement the FFT algorithm. Those that do also gen-

erally provide a routine for calculating the discrete Fourier (spatial) frequencies kl that

match the corresponding Fl ; that is, the kl start from 0 and go to (N/2 1) k and then

from (N/2) k to k. Most FFT packages also have routines that reorder the frequencies

from lowest (most negative) to highest (most positive), and, respectively, the corresponding

Fourier components.2

Once you have obtained the DFT from an FFT algorithm or by direct computation, you

can make direct connection with the conventional continuous Fourier transform and its

inverse simply by supplying x as follows:

0N 1 1

BBBX CCC

F(l k) ! Fl x = BB@ B fn e i (2/N)ln CCC x (7.103)

A

n=0

0 N/2 1 1

1 BBB 1 X CCC 1

f (n x) ! fn = BB@B Fl e i (2/N)ln CCAC . (7.104)

x N l= N/2 x

|x x0 |/a

f (x) = e cos k0 x (7.105)

1

and N = 512 so that x = L/N = 64 . The wave packet given by Eq. (7.105) is plotted in Fig.

7.9 with the points at which fn = f (n x) indicated. We wish to find the DFT for the points

given by fn = f (n x). We use the FFT algorithm to calculate the DFT. The algorithm works

2 For example, in Python, the DFT of an array y is provided within the scipy.fftpack package as fft(y).

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199 Discrete Fourier Transforms

f (x) or f n

0.8 (a)

0.4

0.0 x

1 2 3 4 5 6 7 8

0.4

0.8

Fl

40 (b)

20

0 l

128 256 384 512

20

40

40 (c)

20

0 l

240 256 272 288

20

40

F(k)

(d)

0.4

t

k

20 10 10 20

0.4

Fig. 7.9 Discrete Fourier transform of an exponential wave packet. (a) Exponential wave

packet given by Eq. (7.105). (b) Discrete Fourier transform vs. l. (c) Discrete Fourier

transform reordered so that negative spatial frequencies (l < 256) appear before the

positive spatial frequencies (l > 256). Dotted lines connect DFT points for clarity of

presentation. (d) The DFT multiplied by x is plotted against the proper spatial

frequencies. The continuous solid line shows the analytical Fourier transform.

best if the number of points N in fn is a power of two. That is why we chose N = 512 = 29 .

The DFT is plotted in Fig. 7.9(b) with the positive spatial frequencies running from 0 to 255

and the negative spatial frequencies running from 256 to 511. The DFT is replotted in Fig.

7.9(c) but this time with the spatial frequencies increasing monotonically from negative

(n = 0 to 255) to positive (n = 256 to 511) spatial frequencies.3

The raw DFT is simply indexed by the (dimensionless) integers l, but often you want to

work with the DFT as a function of the proper spatial frequencies, with the proper units of

inverse length. This is a simple matter, you simply shift the indices l by N/2 so that l runs

from N/2 to N/2 1 and then multiply by k. Before performing this rescaling, you must

be sure that you have reordered the DFT as described in the previous paragraph so that the

3 As do many other computer languages, Python provides an FFT helper routine fftshift that shifts the order

of the FFT so that the negative frequencies come first, followed by the positive frequencies.

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200 Fourier Analysis

DFT is plotted starting with the negative spatial frequencies and ending with the positive

spatial frequencies.

Moreover, if you want the digital Fourier transform itself to have the same units as the

equivalent continuous Fourier transform would have, you must multiply the DFT by x,

the spacing between the point of the original data, as indicated by Eq. (7.103). Figure

7.9(d) shows the DFT multiplied by x and plotted against the proper dimensional spatial

frequencies k. For comparison, Fig. 7.9(d) also shows the analytic Fourier transform of

f (x), which is given by

!

a a

F(k) = + eik0 x , (7.106)

1 + (k + k0 )2 a2 1 + (k k0 )2 a2

where we have used Eq. (7.93) to calculate the Fourier transform. Note that because the

original function f (x) is centered about a finite value x0 = 4, there are oscillations, given

by the the exponential factor eik0 x . As expected, the DFT and the analytic Fourier transform

are consistent with each other.

Fourier analysis is extremely powerful and can be used in clever ways that may not

be immediately obvious. One of the more useful applications of Fourier transforms is for

filtering data. Lets look at one illustrative example. Consider the noisy Gaussian waveform

f (x) shown in in Fig. 7.10(a). The DFT is shown in Fig. 7.10(b). What is not so apparent in

Fig. 7.10(b) is that there is some high spatial frequency components in the DFT due to the

noise in f (x). The DFT is shown with a magnified y-axis in Fig. 7.10(c), where the high

frequency noise is clearly visible. The noise in the high frequency components visible in

Fig. 7.10(c) and indicated by the gray arrows is removed in Fig. 7.10(d). Performing the

inverse DFT on the signal in in Fig. 7.10(d) returns the original Gaussian signal with the

noise removed.

You may have noticed that the DFT in Fig. 7.10(b) has oscillations in it. These arise

because the DFT assumes that f (x) goes from 0 to L, which means that all the x values are

positive. The DFT is eectively calculating the Fourier transform of f (x x0 ) where x0 is

half the width of the data: x0 = 4. According to frequency shifting theorem Eq. (7.87), the

1 2 2

resultant Fourier transform is F(k)e ikx0 , where in this case F(k) = e 2 k a . Thus, the factor

of e ikx0 introduces oscillations in the DFT. We could manually remove these oscillations

by multiplying the DFT by eikx0 . However, there is no need to do so as our purpose here is

only to filter out the high frequency components and then re-transform in order to recover

the original signal, less the noise. Once we obtain the Fourier filtered data, which is just a

list of fn values, we are free to plot those values against the original x data, which is what

we do in Fig. 7.10(d).

Finally we point out that the DFT is in general complex, with real and imaginary com-

ponents. In the examples above, the imaginary components are much smaller than the real

components. Nevertheless the imaginary components can be important, especially when

performing a DFT, manipulating it as we have done in the Fourier filtering example above,

and then performing the inverse DFT.

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201 Problems

f (x)

(a)

1.0

0.5

x

4 3 2 1 0 1 2 3 4

F

100 l

(b)

50

0 l

64 128 192 256

50

F

2 l

(c)

0 l

64 128 192 256

F

2 l

(d)

0 l

64 128 192 256

f (x)

(e)

1.0

0.5

t

x

4 3 2 1 0 1 2 3 4

Fig. 7.10 Filtering a noisy Gaussian pulse. (a) Gaussian pulse with noise. (b) DFT of Gaussian

pulse with spatial frequencies shifted so that they appear from negative to positive. (c)

DFT magnified along y direction so that noise is visible. Gray arrows indicate high

frequency noise. (d) Filtered DFT (magnified) with high frequencies removed (e)

Inverse DFT of filtered DFT plotted against x rather than DFT index.

Problems

7.1 Find the Fourier series for the following functions. Take note of the odd or even

symmetry of the functions to avoid calculating Fourier coefficients that must be zero

by symmetry. Sketch the function in each case.

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202 Fourier Analysis

(a)

(

x , if 0x1

f (x) = (7.107)

x , if 1x0.

(b)

(

0 , if 0x1

f (x) = (7.108)

x , if 1x0.

(c) Taking k0 to be a constant, what is the lowest spatial frequency in the Fourier

series for this function and how does it compare to k0 ?

(

sin k0 x , if k0 x 0

f (x) = (7.109)

sin k0 x , if 0 k0 x .

7.2 The following periodic function is defined on the interval [ 2, 2] but repeats indefi-

nitely along the x axis:

8

>

>

>

> 6(2 + x) , if 2 x 32

<

f (x) = >

> 2x , if |x| 32 (7.110)

>

>

: 6(2 x) , if 3

2 x 2 .

(a) Explain how the Python function below codes for f (x). In particular, explain pre-

cisely how the mod works. Then use the function to plot the function over the in-

terval 5 < x < 5. Hint: look at the output of the expression x = np.mod(x+2.0,

4.0) - 2.0 for dierent values of x over the interval 5 < x < 5.

import numpy as np

def sawtooth(x):

x = np.mod(x+2.0, 4.0) - 2.0

y = np.where(np.abs(x) <= 1.5, 2.*x,

6.*(2.-np.abs(x))*np.sign(x))

return y

(b) Find the Fourier coefficients for the sine and cosine Fourier series. You will need

to perform the Fourier integral in a piecewise fashion splitting it up into dierent

sections over the integration interval.

On separate plots, use Python to plot the Fourier series keeping (i) the first non-

zero term in the Fourier series, (ii) the first 2 non-zero terms in the Fourier series,

(iii) the first 3 non-zero terms in the Fourier series, and (iv) the first 10 non-zero

terms in the Fourier series.

(c) Find the power spectrum of f (x), that is the power P(kn ) at each wavevector kn

and make a plot of P(kn ) vs. kn . Find the total power in the waveform by summing

the contributions from all Fourier components. Which Fourier component carries

the most energy?

(d) Find the Fourier coefficients for the exponential Fourier series using Eq. (7.39).

By including both positive and negative coefficients, show that the complex

Fourier series gives the same results as the sine and cosine series for the first

three non-zero terms in the Fourier Series.

t

203 Problems

7.3 Find the Fourier transforms of the following functions. Using Python, plot the func-

tions below and on a separate graph, plot both the real and imaginary parts (if they

exist) of the Fourier transforms of each function. From your plots, estimate k (or

!) and x (or t) and show that k x 1 (or ! t 1).

Take care about the absolute values in the definitions of f1 (x) and f2 (x) below

when performing the Fourier integrals.

(a) f1 (x) = 5e 3|x|

(b) f2 (x) =(5e 3|x| sin 2x

2(1 14 t2 ) , if |t| 2

(c) g(t) =

0 , if |t| > 2 .

7.4 The (spatial) frequency shifting theorem says the if f (x) and F(k) are Fourier trans-

form pairs, i.e. f (x) () F(k) [sometimes written as F { f (x)} = F(k)], then

f (x) eik0 x () F(k k0 ) (7.111)

(a) Prove the frequency shifting theorem. That is, given that F(k) is the Fourier

transform of f (x), show that Eq. (7.111) is true.

(b) In Exercises 7.3(a) and 7.3(b), you found the Fourier transforms of the functions

f1 (x) = 5e 3|x| and f2 (x) = 5e 3|x| sin 2x by explicitly calculating the Fourier inte-

grals for each function. For this exercise, you are to obtain the Fourier transform

of f2 (x) = 5e 3|x| sin 2x by applying the frequency shifting theorem to the Fourier

transform of f1 (x) = 5e 3|x| . This illustrates how the frequency shifting theorem

can make the calculation of the Fourier transforms of functions involving sines

and cosines significantly easier.

7.5 Use the FFT feature of SciPy to find the discrete Fourier transforms of the functions

listed in problem 7.3. Take care to adequately sample, but not drastically oversample

the three functions. The number of points you use should be a power of 2. Plot the

discrete Fourier transforms you calculate and compare them (properly shifted and

scaled) to the continuous Fourier transforms you determined analytically in problem

7.3. After finding the discrete Fourier transforms, perform the inverse transforms to

and show that you get back the original function in each case.

7.6 This problem explores the use of spatial filtering with FFTs to clean up noisy data.

Start by considering the following Python routine. It generates a noisy Guassian

function and then performs a discrete Fourier transform (DFT). It then zeros the

high frequency components (recall the nonintuitive ordering of the frequenies for

FFT routines), and then inverse transforms the filtered transform to return a clean

Gaussian signal.

from numpy.fft import fft, fftshift, ifft

import numpy as np

import matplotlib.pyplot as plt

return np.exp(-x*x/2.0) * (1.0 + noiseAmp*(np.random.randn(len(x))))

t

204 Fourier Analysis

N = 256

x = np.linspace(-4.0, 4.0, N)

y = gaussNoisy(x, 0.1)

YFT = fftshift(yFT) # DFT reordered so that - frequencies are first

plt.subplot(511)

plt.plot(x,y) # Noisy Gaussian function

plt.subplot(512)

plt.plot(YFT.imag) # Imaginary part of DFT

plt.subplot(513)

plt.plot(YFT.real) # Real part of DFT

plt.subplot(514)

plt.plot(fftshift(yFT)) # Cleaned up DFT

plt.subplot(515)

plt.plot(yFTFT) # Inverse of cleaned up DFT

plt.show()

(a) Set the noise amplitude to zero in the and find the DFT for the clean Gaussian

function. How far away from the maximum in the DFT do you have to go for all

the DFT amplitudes to fall below 1% of the maximum.

(b) Generate the proper Fourier transform variable k for the Fourier transformed data

and replot F(k) vs. k following the example of Fig. 7.9(d).

(c) Rescale the y-axis for the plot of the DFT so that the high frequency Fourier

components due to the noise are clearly visible, much as we did in Fig. 7.9(c).

(d) In the program above, all the DFT components more than 6 points from the

center spatial frequency are zeroed. What value of k does this correspond to? Is

this a good choice or is there a better choice? Justify your answer.

(e) Why is the imaginary component of the FFT so much smaller than the real com-

ponent? Is there some signal there or is it all noise?

(f) Vary the amplitude of the noise set by the noiseAmp variable. Currently the

t

205 Problems

noise amplitude is set to 0.3 or 30% of the signal. How well does the Fourier

filtering work with the noise is significantly larger, say 100%?

8 Sound and Electromagnetic Waves

8.1 Sound

206

t

207 Problems

Problems

8.1

Interference, Diffraction, and Fourier

9

Optics

9.1 Interference

9.2 Diffraction

9.4 Scattering

208

t

209 Problems

Problems

9.1

PART II

CLOSING FEATURES

A Appendix A Linear algebra

In this Appendix, we review the properties of matrices that are relevant to our study of

coupled oscillators. For our purposes, matrices can be regarded as generalizations of vec-

tors, so before summarizing the properties of matrices, we briefly review vectors. We then

provide a concise overview of those aspects of linear algebra and matrices that are relevant

to our study of oscillations and normal modes.

A.1 Vectors

A vector is an ordered list of numbers. The numbers that make up a vector are called its

components. The number of components in the list is called the dimension of the vector. For

example, we can represent the spatial location of an point in space by a three-dimensional

vector, which we might express in terms of its Cartesian components: r = (r x , ry , rz ) =

(x, y, z). In general, higher dimensional vectors can be constructed, say with N components:

S = (s1 , s2 , s3 , ..., sN ). Vectors also obey a number of familiar rules pertaining to how

they are added, subtracted, and multiplied. Here we summarize those rules for the vectors

a = (a1 , a2 , a3 ) and b = (b1 , b2 , b3 ):

a b = (a1 b1 , a2 b2 , a3 b3 ) dot product (A.2)

a b = (a2 b3 a3 b2 , a3 b1 a1 b3 , a1 b2 a2 b1 ) cross product (A.3)

Scalar multiplication and the dot product are commutative. That is, the order of multipli-

cation doesnt matter so that ta = at and a b = b a. By contrast, the cross product does

not commute; in general a b , b a, as can be readily demonstrated using the above

definition.

Each of the three types of vector multiplication have a geometrical interpretation. Scalar

multiplication of a by t means that the vector a is stretched by a factor or t. The dot product

a b represents a scalar with a length equal to the projection of a onto b or of b onto a. The

cross product a b gives a vector representing the rotation of a into b.

213

t

214 Linear algebra

A.2 Matrices

Matrices are arrays of numbers arranged in rows and columns. A 2 row by 3 column matrix,

or 2 3 matrix, A is defined as

!

A11 A12 A13

A= . (A.4)

A21 A22 A23

We denote matrices by bold sans-serif characters and their elements by the same sub-

scripted italic characters. We shall often work with square matrices, meaning matrices that

have the same number of rows and columns

0 1

BBBA11 A12 A13 CCC

B C

A = BBBBA21 A22 A23 CCCC . (A.5)

B@ CA

A31 A32 A33

rows, i.e. ATij = A ji . Thus, the transpose of a matrix with 2 rows and 3 columns would have

3 rows and 2 columns.

With these definitions, we see that vectors are simply matrices that have either one row or

one column. When working with matrices, we generally write a column vector (1 column,

n rows) as a and a row vector (n columns, 1 row) as the transpose of a column vector a T .

Matrices can be added and subtracted. Like vectors, they can also be multiplied in several

dierent ways.

Matrix addition

Matrix addition and subtraction are defined on an element-by-element basis. Therefore, to

be added or subtracted, two matrices must have the same number or rows and columns. For

the 2 3 matrices A and B , this means that

! !

A11 A12 A13 B B12 B13

C = A +B = + 11 (A.6)

A21 A22 A23 B21 B22 B23

!

A11 + B11 A12 + B12 A13 + B13

= , (A.7)

A21 + B21 A22 + B22 A23 + B23

subtraction for C = A B is defined in the obvious way Ci j = Ai j Bi j .

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215 Properties of matrices

Matrix multiplication

As in the case of vectors, several types of multiplication are defined for matrices. In scalar

multiplication, all components of a matrix are multiplied by the same scalar value: the

components of tA , where t is a scalar, are given by t Ai j for all indices i j.

The matrix dot product C = A B , often simply called matrix multiplication and written

C = AB , is defined by the rule

X

Ci j = Aik Bk j , (A.8)

k

Ci j = Aik Bk j , (A.9)

where it is understood that repeated indices (in this case k) are summed over. This definition

only makes sense, and thus the matrix dot product is only defined, if the number of columns

in A is equal to the number of rows in B . The matrix dot product is also often referred to

as the inner product. Except in special cases, matrix multiplication is not commutative,

that is, AB , BA in general. Matrices that do commute must be square, because of the

definition of matrix multiplication. This is a necessary but not sufficient condition for two

matrices to commute.

The matrix outer or direct product is also defined, and is usually designated as A B .

But as we will not need it, we do not discuss it further.

Just as vectors have a number of properties, such as length and direction, matrices also

have a number of properties that can be characterized by various well-defined quantities.

Here we review a few of the more relevant of those properties.

The determinant is defined for square matrices and is designated by absolute value signs

(or in some texts as det A ):

A11 A12 A13 A1n

A21 A22 A23 A2n

|A | = det A = A31 A32 A33 A3n (A.10)

.. .. .. .. ..

. . . . .

An1 An2 An3 Ann

The determinant of a 2 2 matrix is given by

A11 A12

|A | = = A11 A22 A12 A12 (A.11)

A12 A22

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216 Linear algebra

The determinant of any higher order matrix can be calculated using Laplaces formula

n

X

|A| = ( 1)1+ j A1 j M1 j , (A.12)

j=1

where M 1 j is the (n 1)(n 1) matrix formed by removing the first row and the jth column

of A , the so-called 1 j minor of A . Thus for a 3 3 matrix, the determinant is given by

|A | = A21 A22 A23 (A.13)

A31 A32 A33

A22 A23 A A23 A A22

= A11 A12 21 + A13 21 (A.14)

A32 A33 A31 A33 A31 A32

= A11 (A22 A33 A23 A32 ) A12 (A21 A33 A23 A31 + A13 (A21 A32 A22 A31 ) .(A.15)

Larger matrices can be broken down in terms of their minors in a similar fashion. In fact,

there is no need to use first row as the special row in the definition of |A | as we did in

Eq. (A.12). Any row or column may be designated as the special" row or column around

which the determinant is calculated, but it is often convenient to choose the first row, as we

have done above.

Exercise A.3.1 Using the definition of the determinant given by Eq. (A.12), show that

the the determinant of a lower or upper triangular matrix is given by the produce

P P

of its diagonal elements. That is, show that |L | = ni=1 Lii and that |U | = ni=1 Uii ,

where L is a lower triangular matrix, meaning that all the elements above the matrix

diagonal are zero, and where U is a upper triangular matrix, meaning that all the

elements below the matrix diagonal are zero.

1

The inverse of a matrix A is designated as A and has the property that A 1A = I where I

is the identity matrix.

A matrix A is said to be symmetric if its elements Ai j obey the relation Ai j = A ji . Thus, for a

symmetric matrix A T = A . Similarly, a matrix B is said to be antisymmetric if its elements

Bi j obey the relation Bi j = B ji . Thus, for an antisymmetric matrix B = B T . Clearly, for

these relations to hold, the matrices A and B must each have the same numbers of rows and

columns, so only square matrices can be symmetric or antisymmetric.

Of course, most square matrices are neither symmetric or antisymmetric. However, its

easy to show that any square matrix can be expressed as the unique sum of a symmetric

and antisymmetric matrix. See if you can show that.

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217 Properties of matrices

A matrix A is said to be Hermitian1 if the complex conjugate of the transpose of the matrix

is equal to the matrix itself, i.e. if A T = A . The complex conjugate of the transpose of

the matrix A is called the adjoint and is often denoted as A , i.e. A A T . Therefore,

Hermitian matrices are matrices for which A = A and are thus said to be self-adjoint.

An n n square matrix P is said to be positive definite if a T P a > 0 for all vectors a with n

real components. If a contains complex components, then P is said to be positive definite

if aP a > 0 for all vectors a with n complex components, where a is the transpose of

the complex conjugate of a. Clearly, the first definition is a special case of the second. It

is straightforward to show that all Hermition matrices are positive definite. Because a real

symmetric matrix is Hermitian, real symmetric matrices are also positive definite.

Positive definite matrices have a number of useful properties. Below we list those that

are particularly useful for our purposes.

The eigenvalues of positive definite matrices are all real and positive.

Positive definite matrices P can always be factorized such that P = LL where L is a

unique lower triangular matrix with diagonal elements that are all positive. This factor-

ization is called Cholesky decomposition. If P is real, then P = LL T where L is real.

Two positive definite matrices can be simultaneously diagonalized, that is, diagonalzed

by the same set of similarity transformations.

course this means that for an orthogonal matrix AA T = I . The determinant of an orthogonal

matrix is either 1 or -1; that is, |A | = 1.

We list here a few useful matrix identities, most of which are easily proved:

1

(AB ) = B 1A 1

(AB )T = B T A T (A 1 )T = (A T ) 1

|AB | = |A | |B | |A 1 | = 1/|A |

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218 Linear algebra

A.4 Transformations

Matrices that operate on (e.g. multiply) vectors or other matrices represent linear transfor-

mations. The kinds of transformations represented by matrix operations can be classified

into a number of dierent types, many of which have important geometrical or mathemat-

ical properties. Below we consider a few that are useful for our studies.

A.4.1 Rotations

In linear algebra it is often useful to write a square matrix A as the product of a lower

triangular matrix L and an upper triangular matrix U . That is,

A = LU (A.16)

(A.17)

or

0 1 0 10 1

BBBA11 A12 A1n CCC BBBL11 0 0 CCC BBBU11 U12 U1n CCC

BBBBA A22

C B

A2n CCCC BBBBL21 L22 0 CCCC BBBB 0 U22

C

U2n CCCC

BBB 21 C B CCC BBB C

BBB . .. .. .. CCCC = BBBB .. .. .. .. CCC BBB . .. .. .. CCCC .

BBB .. . . . CCCC BBBB . . . . CCC BBB .. . . . CCCC

B@ A @ CA B@ A

An1 An2 Ann Ln1 Ln2 Lnn 0 0 Unn

Factoring a matrix this way is known as LU decomposition and is generally useful in nu-

merical methods for solving systems of linear equations. Such a factorization is not unique;

there are many LU decompositions for a given matrix. For square symmetric matrices S

that are positive definite (A.3.5), it can shown that the LU decomposition can be done

such that U = L T , that is

S = LL T . (A.18)

Lets start by explicitly calculating L for a general 3 3 matrix:

S = LL T (A.19)

or

0 1 0 10 1

BBBS 11 S 21 S 31 CCC BBBL11 0 0 CCC BBBL11 L21 L31 CCC

BBB CCC BBB CCC BBB C

BBBS 21 S 22 S 32 CCC = BBBL21 L22 0 CCC BBB 0 L22 L32 CCCCC (A.20)

@ A @ A@ A

S 31 S 32 S 33 L31 L32 L33 0 0 L33

0 2 1

BBB L11 L11 L21 L11 L31 CCC

BBB C

= BBBL11 L21 2

L21 + L22 2

L21 L31 + L22 L32 CCCCC . (A.21)

@ 2 2 2 A

L11 L31 L21 L31 + L22 L32 L31 + L32 + L33

t

219 Transformations

Note that the matrix LL T is symmetric, as it should be. We can determine the elements of

L by equating S element by element with LL T starting with the first column. Thus we have

2

p

S 11 = L11 ) L11 = S 11 (A.22)

S 21 = L11 L21 ) L21 = S 21 /L11 (A.23)

S 31 = L11 L31 ) L31 = S 31 /L11 (A.24)

Notice that Eqs. (A.23) and (A.24) are obtained using L11 the result of Eq. (A.22). Contin-

uing with the second column, starting with the diagonal element we obtain

q

2 2 2

S 22 = L21 + L22 ) L22 = S 22 L21 (A.25)

S 32 = L21 L31 + L22 L32 ) L32 = (S 32 L21 L31 )/L22 . (A.26)

Once again we see that Eqs. (A.25) and (A.26) both use results from previous calculations

of various matrix elements of L . Finally, from the third column we obtain

q

2 2 2 2 2

S 33 = L31 + L32 + L33 ) L33 = S 33 L31 L32 . (A.27)

Generalizing these results to an N N matrix, the diagonal terms are given by

v

u

t j 1

X

Ljj = S jj L2jk . (A.28)

k=1

0 j 1

1

1 BBBB X CCC

Li j = B@BS i j Lik L jk CCCA . (A.29)

Ljj k=1

Equations (A.28) and (A.29) are suitable for using in a computer routine to numerically

compute the matrix L for any given S . Our assumption that S is a positive definite matrix

guarantees that the argument of the square root in Eq. (A.28) is always positive. Here is a

Python function that implements Eqs. (A.28) and (A.29) to calculate L :

import numpy as np

def choleski(s):

n = len(s)

for k in range(n):

try:

s[k,k] = np.sqrt(s[k,k] - np.dot(s[k,0:k],s[k,0:k]))

except ValueError: # for square root of negative number

print("Matrix is not positive definite")

for i in range(k+1,n):

s[i,k] = (s[i,k] - np.dot(s[i,0:k],s[k,0:k]))/s[k,k]

for k in range(1,n): # put zeros in upper triangle

s[0:k,k] = 0.0

return s

t

220 Linear algebra

The program performs the Choleski decomposition in place" meaning that it puts the

result into the same matrix that was input into the routine. This has the virtue of being

faster and saving computer memory. However, if you want to save the original matrix,

you had better make a copy of it using the NumPy copy function before performing the

Choleski decomposition.

Exercise A.4.1 Perform Choleski decomposition of the following symmetric matrix

by hand using the algorithm described by Eqs. (A.28) and (A.29). Verify that the

algorithm works in this case by explicitly showing that S = LL T .

0 1

BBB 6 20 15CCC

BBB C

S = BB20 80 50CCCC

B@ CA

15 50 60

Exercise A.4.2 Use the Python function given above to find the Choleski decompo-

sition of the matrix S above. Verify that you get the same answer that you got in

Problem 1 above. Explain how the program works, line by line. Include in your ex-

planation why it is not a problem that this routine calculates L in place.

The solution to a system of homogeneous linear equations of the form C a = 0 exists only

if det C = 0. We can easily prove this for a general 2 2 matrix C with matrix elements Ci j

and column vector a with elements a1 and a2 by writing out the matrix equation C a = 0

explicitly

! !

C11 C12 a1

=0. (A.30)

C21 C22 a2

or

C11 a1 + C12 a2 = 0 (A.31)

C21 a1 + C22 a2 = 0 (A.32)

From Eq. (A.32),

C21

a2 = a1 . (A.33)

C22

Substituting this into Eq. (A.31) yields

C21

C11 a1 C12 a1 = 0 (A.34)

C22

(C11C22 C12C21 ) a1 = 0 (A.35)

det C (C11C22 C12C21 ) = 0 . (A.36)

t

221 Eigenvalue equations

uninteresting case in which a1 = a2 = 0. If C21 = 0, the C22 a2 = 0 and either C22 = 0 or

a2 = 0. The latter case is uninteresting again with a1 = a2 = 0. If C22 = 0 and C21 = 0, then

det C = 0, which is uninteresting. The end result is that either a1 = a2 = 0 or det C = 0.

Lagrangian formulation of classical

B

mechanics

You are used to finding the equation of motion in classical mechanics by careful application

Newtons second law, F = ma. But there is more than one way to skin a cat" as the

saying goes. In this appendix, we show you a versatile and powerful method for finding

the equation of motion for a dynamical system.

The method was developed in the 18th century by Joseph Louis Lagrange and constitutes

an elegant and powerful reformulation of Newtonian mechanics. The method is particularly

useful when there are elements in the problem that constrain the motion of the system, as

the method eliminates the need to consider the forces associated with those constraints. An

example is the rod in a pendulum that constrains the mass at the end of the rod to move in

an arc. The Lagrangian formulation eliminates the need to consider the tension in the rod: it

simply never appears in the formulation or solution of the problem. This approach becomes

particularly useful in more complex problems, such as the double pendulum presented

below, in which the constraint forces can be quite a nuisance to determine.

B.1 Lagrangian

The Lagrangian formulation of classical mechanics begins with a quantity called the La-

grangian L, which is the dierence the kinetic energy K and potential energy U of a system.

L=K U. (B.1)

For a simple unconstrained particle, the kinetic energy in Cartesian coordinates is given by

1 2

K= m x + y2 + z2 . (B.2)

2

As the potential energy is a function of the coordinates x, y, and z, the Lagrangian is thus a

function of the coordinates and velocities of the particle: L = L(x, y, z, x, y, z). The positions

(x, y, z) and velocities ( x, y, z) are the independent variables of the Lagrangian. The force

on the particle is given by the gradient of L, or in component form

@L @U

= = Fx , (B.3)

@x @x

222

t

223 Lagrangian

with similar equations for Fy and Fz . Taking derivatives of L with respect to the velocities

gives the various components of the momentum

@L

= m x = p x , (B.4)

@ x

with similar equations for py and pz . Noting that Newtons second law is most generally

given by F = dp/dt we can use these results to rewrite Newtons second lawn terms of the

Lagrangian function L as

@L d @L @L d @L @L d @L

= , = , = . (B.5)

@x dt @ x @y dt @y @z dt @z

Thus far, in writing down Eq. (B.5) we have just managed to rewrite F = ma in a somewhat

more complicated form. We have yet to see the utility or power of this formulation of

Newtons second law. The power of this formulation only becomes apparent when we

consider constrained systems whose motion is described by generalized coordinates. Dont

expect that last sentence to make much sense yet. Just be patient and follow the argument

laid out below.

Suppose we wish to describe the state of our system using some set of coordinates other

than Cartesian coordinates. For example, we may wish to use angular coordinates, as we

do for the pendulum in Chapters 13. For more complex problems, it is often convenient

to use a mixture of translational and orientational coordinates. For such cases, it would be

useful to reformulate Eq. (B.5) in terms of generalized coordinates.

As a first step, let us recognize that it will take a certain minimal number of coordinates

to specify the configuration of a system. For the simple pendulum described in Chapters

13, that number is 1, say the angle or the arc length s, as defined in Fig. 1.1. Of course,

we can also describe the pendulums position in terms of x and y, but we can eliminate the

y coordinate in favor x using the geometry specified by the constraint that the pendulum

move along an arc of constant radius l, the length of the pendulum. In other cases, say

for the coupled double pendulum discussed in 4.1.1, two coordinates may be required. In

other cases, still more coordinates may be required, but in general some minimum number

of generalized coordinates will be required. Lets designate those generalized coordinates

as q1 , q2 , ..., qN . If each coordinate qi can be varied independently of all the other coordi-

nates q j , qi , then the coordinates are said to be holonomic. In this case, the number of

coordinates N is equal to the number of degrees of freedom of the system.

In some cases, the system of coordinates do not all vary independently. The classic ex-

ample is that of a ball rolling without slipping on a surface. In such a case, two coordinates

are required to specify translation of the center of mass of the sphere and three more (an-

gular coordinates) are required to specify its orientation. However, if the ball moves, at

least two coordinates, one for translation and one for orientation change. Moreover, there

is no simple way to eliminate one coordinate in favor of another in this case. Such coordi-

nates, which cannot all vary independently, are said to be nonholonomic. Systems requiring

nonholonomic coordinates are generally much more difficult to treat. Fortunately, we shall

need to consider only systems that can be described with holonomic coordinates. In what

follows, therefore, we assume that we can find coordinates that are holonomic.

With this in mind, we can specify the state of a system with N degrees of freedom in

t

224 Lagrangian formulation of classical mechanics

terms of N Cartesian coordinates {ri }, which would encompass the x, y, and z coordinates

of all the particles in the system. The kinetic energy of such a system is given by

N

X 1

K= mi ri2 , (B.6)

i

2

where the sum extends over all N degrees of freedom. In order to generalize Eq. (B.5),

we would like to express the kinetic energy in terms of the generalized coordinates {qk }.

In general, we can express each of the N Cartesian coordinates ri as functions of the N

generalized coordinates {qk }:

where we have also included an explicit time dependence for the case the dependence of

the Cartesian and generalized coordinates is time dependent (in some known way). Thus,

X @ri @qk @ri

ri = + (B.8)

k

@qk @t @t

X @ri @ri

= qk + (B.9)

k

@qk @t

Taking the derivative @ri /@qm picks out of the sum in Eq. (B.9) only that term where k = m,

as all other derivatives are zero, yielding

@ri @ri

= . (B.10)

@qm @qm

Of course, the index label is arbitrary so we can also write the above equation with m = k.

Multiplying both sides by r and taking the time derivative gives

! !

d @ri d @ri

ri = ri (B.11)

dt @qk dt @qk

@ri @ri

= ri + ri . (B.12)

@qk @qk

As it doesnt matter in which order we perform the dierentiations, we can rewrite this as

0 1 0 1

d BBB @ ri2 CCC @ri @ BBB ri2 CCC

@B AC = ri + @B AC . (B.13)

dt @qk 2 @qk @qk 2

Multiplying both sides by mi gives

0 1 0 21

d BBB @ mi ri2 CCC @r @ BBB mi ri CCC

B@ CA = mi ri i + B@ CA . (B.14)

dt @qk 2 @qk @qk 2

Summing over i yields

! X

d @K @ri @K

= Fi + , (B.15)

dt @qk i

@qk @qk

t

225 Lagrangian

where we have used Newtons second law Fi = mi ri and Eq. (B.6) for the kinetic energy.

Rearranging terms, we can rewrite this equation as

!

d @K @K

= Fk , (B.16)

dt @qk @qk

where Fk is known as a generalized force, which is defined as

X @ri

Fk Fi . (B.17)

i

@qk

Equation (B.16) is completely general so long as the generalized coordinates {qk } are all

independent. A further simplification is possible if the force is conservative, meaning that

it can be written as the gradient of a potential energy, or in component form, Fi = @U/@ri .

In this case the generalized force can be written as

X @ri X @U @ri

Fk = Fi = (B.18)

i

@qk i

@ri @qk

Equation (B.7) tells us that we can write each Cartesian coordinate ri in terms of general-

ized coordinates, which in turn tells us that we can write the potential energy as a function

of the generalized coordinates U(q1 , q2 , ..., qN ). Taking the derivative with respect to qk

gives

@U X @U @ri

= . (B.19)

@qk i

@ri @qk

Thus we see that the generalized force for a conservative system is

@U

Fk = . (B.20)

@qk

Substituting this expression into Eq. (B.16) gives

!

d @K @K @U

= . (B.21)

dt @qk @qk @qk

Finally, using our definition of the Lagrangian L = K U, we obtain our final result

!

d @L @L

= , (B.22)

dt @qk @qk

where L is understood to be a function of the generalized coordinates {qi , qi } and we have

used the fact that U depends only on the coordinates {qi } so that @U/@qi = 0. Equation

(B.22) expresses the Lagrangian reformulation of Newtons second law for conservative

systems, that is, for systems without any dissipation (e.g. friction). Because it is expressed

in generalized coordinates, it can be significantly simpler to apply to a dynamical system

than a straightforward application of Newtons second law. In the next section we illustrate

its utility with a couple of examples.

If there is dissipation in the problem, then Eq. (B.16) must be used. In fact, one can

replace K by L in Eq. (B.16) for all the conservative forces in the problem and then only

calculate the generalized force F for the nonconservative forces. We will not pursue further

the application of Lagranges equations to systems with nonconservative forces.

t

226 Lagrangian formulation of classical mechanics

Before working some examples, we wish to make one final point. Just as Eq. (B.18)

defines a generalized force for a conservative system, we can also define a generalized

momentum

@L

Pk = = mqk . (B.23)

@qk

In the examples that follow, we see that the generalized forces and momenta turn out to be

actual forces and momenta or actual torques and angular momenta, depending on whether

the generalized coordinates associated with them are translational or orientational coordi-

nates.

Lets see how these results work a few practical examples. We begin by considering an

almost trivial example, the undamped pendulum pictured in Fig. 1.1. There is only one

degree of freedom in this problem so we need only one generalized coordinate, which we

take to be . The kinetic and potential energies, expressed in terms of this generalized

coordinate, are

1 2 2

K= ml , U = mgl(1 cos ) , (B.24)

2

where we have taken U = 0 to be where = 0. Therefore

1

L = K U = ml2 2 mgl(1 cos ) . (B.25)

2

The derivatives in Eq. (B.22) are

@L @L

= ml2 , = mgl sin . (B.26)

@ @

Thus, taking time derivative @L/@ = ml2 ,

Eq. (B.22) gives

which is the equation of motion for the simple pendulum (cf. Eq. (1.3)).

Notice that we obtained the equation of motion, an equation that involves forces, en-

tirely from energy considerations. No forces were used. Moreover, there was nowhere any

reference to the constraint force provided by the tension in the (massless) pendulum rod.

In this case, the generalized force is

@U

F = = mgl sin , (B.28)

@

which is a torque, as one might expect since the generalized coordinate is an angle. Simi-

larly, the generalized momentum is actually the angular momentum

@L

P = = ml2 . (B.29)

@

t

227 Double pendulum

To see the real power of the Lagrangian formulation, we turn to a more complex problem:

the double pendulum. This problem is deceptively difficult if one tries to apply Newtons

second law in a straightforward way. As we shall see, it is quite manageable when analyzed

using the Lagrangian formulation.

The double pendulum is shown in Fig. B.1. For simplicity we assume that the motion of

the pendulum is confined to a single plane. The kinetic energy of the upper mass is 12 ml2 2 .

The kinetic energy of the lower mass is quite a bit more complicated, as its velocity is the

as indicated in Fig. B.1. The resultant velocity

vector sum of the two velocities l and l,

can be obtained from the law of cosines where the angle between the two vectors is .

The total kinetic energy is thus given by

1 2 2 1 2 2

K= ml + ml + 2 cos( ) + 2 . (B.30)

2 2

The potential energy of the system is of the upper mass is mgl(1 cos ) and of the lower

pendulum is mgl[(1 cos ) + (1 cos )]. The total potential energy of the system is thus

U = mgl 2(1 cos ) + (1 cos ) . (B.31)

L=K U

!

1

= ml 2 + cos(

2

) + 2 mgl 2(1 cos ) + (1 cos ) . (B.32)

2

t

Fig. B.1 Double pendulum.

t

228 Lagrangian formulation of classical mechanics

It is straightforward, albeit a bit tedious, to apply Eq. (B.22), Lagranges equations, to ob-

tain the equations of motion for and . The result is a pair of highly nonlinear equations

that cannot be solved analytically. To obtain an analytical solution, you have to make the

usual simplifying small-angle assumption in order to linearize the equations of motion.

Knowing this in advance, it is simpler to make the small angle approximation on the La-

grangian itself. But, as Lagranges equations involve derivatives, it is necessary to keep

terms up to quadratic order in and (and and ). Thus, we make the approximations

cos ' 1 12 2 and cos ' 1 12 2 . We make the approximation cos( ) ' 1 as

cos(

) is multiplied by , which is already quadratic in order, meaning that keeping

any higher order terms in the Taylor series expansion of cos( ) would lead the third and

higher order terms which would be discarded in the small angle approximation. Making

these small angle approximations, Eq. (B.32) becomes

L = ml2 2 + + 12 2 mgl 2 + 12 2 . (B.33)

Lagranges equations of motion for this system are

! !

d @L @L d @L @L

= , = , (B.34)

dt @ @ dt @ @

which become

ml2 (2 + ) = 2mgl (B.35)

ml2 ( + ) = mgl . (B.36)

These equations can be rewritten in matrix form m x = k x where

! ! !

2 2 1 2 0

m = ml , k = mgl , x= . (B.37)

1 1 0 1

Using the trial solution x = aei!t , the matrix equation becomes an eigenvalue equation

G a = a where G m 1k . The inverse of m is

!

1 1 1

m 1= 2 (B.38)

ml 1 2

so that

!

1g 2 1

G =m k = (B.39)

l 2 2

To find the eigenvalues, we solve the secular equation

2g/l g/l

|G I| = =0 (B.40)

2g/l 2g/l

!2 g 2

2g

= 2 =0. (B.41)

l l

which, recalling that = !2 , gives the normal frequencies

r

g p

!= 2 2 (B.42)

l

t

229 Symmetry of the mass and stiffness matrices

(G I) a = 0 . (B.43)

p p

For = !2 = (g/l)(2 2) and for = !2 = (g/l)(2 + 2), this gives the eigenvectors

! !

1 1 1 1

a = p p , a = p p . (B.44)

3 2 3 2

The examples of coupled oscillators discussed in Chapter 4 and in B.3 above illustrate the

general method of finding the normal modes of a system. We start by finding the linearized

equations of motion for a system, which are then cast as an eigenvalue problem. Solving the

problem then consists primarily of find the normal frequencies and eigenvectors (or normal

coordinates) of the system. The Lagrangian formulation provides a very general method

for finding the equations of motion based on first writing down the kinetic and potential

energies in terms of their generalized coordinates {qk , qk }. Because we are interested only

in the linearized equations of motion for the normal mode problem, we consider only small

oscillations about the equilibrium positions of all the dierent degrees of freedom.

For the cases of interest here, the potential energy is a function only of the general-

ized coordinates and not their time derivatives (the generalized velocities"). Using vector

notation for the entire set of normal coordinates, q {qk }, we can express the potential

energy of the system for small excursions from equilibrium by a Taylor series about their

equilibrium positions, defined by q = 0

X @U ! 1 X @2 U

!

U(q) = U(0) + qk + qk ql + ... (B.45)

k

@qk 0 2 k,l @qk @ql 0

Since we are expanding U about the equilibrium positions, the first derivatives (@U/@qk )0

are zero. Thus, to second order in the displacements, we can write the potential energy as

1X

U(q) = kkl qk ql , (B.46)

2 k,l

where

!

@2 U

kkl , (B.47)

@qk @ql 0

and we have chosen the zero of potential energy so that U(0) = 0. Note that kkl = klk , as

the order of dierentiation doesnt matter.

t

230 Lagrangian formulation of classical mechanics

The kinetic energy is a function of the generalized velocities q {qk }, and quite gener-

ally have the form

1X

K(q, q) = mkl (q) qk ql , (B.48)

2 k,l

where the coefficients mkl (q) can depend on the generalized coordinates q (see Eq. (B.30)

for an example). For small deviations from equilibrium we can expand mkl (q) in a Taylor

series about q = 0:

X @mkl

mkl (q) = mkl (0) + q s + ... (B.49)

s

@q s

However, the expression for the kinetic energy Eq. (B.48) is second order in the general-

ized velocities as every term in the sum contains a factor qk ql . Retaining any terms in the

expansion for mkl (q) beyond the constant mkl (0) would lead to third and higher order terms.

We therefore keep only the constant term mkl (0). Thus, for small oscillations, we can write

the kinetic energy as

1X

K(q) = mkl qk ql , (B.50)

2 k,l

where in writing mkl , we mean mkl (0), here and hence forward. This equation serves as the

definition of mkl which, similar to the definition of kkl by Eq. (B.47), can be written as

!

@2 K

mkl , (B.51)

@qk @ql 0

where the kinetic energy used in this equation is the form containing terms only up to

quadratic order given by Eq. (B.50). Note that mkl = mlk as the order of dierentiation

doesnt matter.

The Lagrangian of a linearized system of coupled oscillators is thus given by

L(q, q) = K(q) U(q) (B.52)

1X 1X

= mkl qk ql kkl qk ql . (B.53)

2 k,l 2 k,l

We can obtain the generalized equations of motion for small oscillations of a system of

coupled oscillators from Lagranges equations. Because the kinetic energy is a function

only or the generalized velocities q and the potential energy is a function only or the gen-

eralized coordinates q Lagranges equations become

!

d @K @U

= . (B.54)

dt @qk @qk

Using Eqs. (B.48) and (B.50) for K and U gives

X X

mkl qk = kkl qk , (B.55)

k k

m q = k q . (B.56)

t

231 Problems

Here the coordinates q are understood to be generalized coordinates and not the normal

coordinates. Note that the matrices m and k are symmetric owing to their definitions by

Eqs. (B.51) and (B.47). You can use the fact that m and k are symmetric to check that

you have not made a mistake in setting up the equations of motion for a linearized system

of coupled oscillators, whether or not you use the Lagrangian formulation to obtain the

equations of motion.

Problems

B.1 Consider the problem of a pendulum swinging from an oscillating mass worked in

??. Show that the Lagrangian of the system is given by

1 1 1 2

L= (M + m) x2 + ml x cos + ml2 2 mgl(1 cos ) k x . (B.57)

2 2 2

Use Eq. (B.22), Lagranges equations, to find the equations of motion of the system

and show that you get the same results as those obtained directly from Newtons 2nd

law, which are given in ?? by Eqs. (??) and (??).

C Appendix C Computer programs

...

232

Notes

233

t

235 Notes

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