Module MS4407 Perturbation Techniques and Asymptotics PDF

Module MS4407: Perturbation techniques and asymptotics
Prof. Stephen O’Brien.
Texts
Introduction to perturbation methods, Nayfeh.
Perturbation methods in applied mathematics, Kevorkian/Cole.
Perturbation methods, Hinch.
Advanced mathematical methods for scientists and engineers, Bender/Orszag.
Perturbations, Murdock.
Mathematics applied to deterministic problems in the natural sciences, Lin & Segel.
Objectives and Assessment

This course introduces the basic principles of perturbation methods and asymptotic techniques.
Assessment is via participation & midterm test (10%) and end of term examination (90%).
Repeat examination is in the autumn.
Exam papers
Public folders/All Public Folders/Academic/Exam Papers/Maths & Stats
http://www.staff.ul.ie/obriens/index.html
GREEK ALPHABET
For each Greek letter, we illustrate the form of the capital letter and the form of the lower case letter.
In somecases, there is a popular variation of the lower case letter.
Greek Greek English Greek Greek English
letter name equivalent letter name equivalent
A α Alpha a N ν Nu n
B β Beta b Ξ ξ Xi x
Γ γ Gamma g O o Omicron o
∆ δ Delta d Π π ̟ Pi p
E ǫ ε Epsilon e P ρ ̺ Rho r
Z ζ Zeta z Σ σ ς Sigma s
H η Eta e T τ Tau t
Θ θ ϑ Theta th Υ υ Upsilon u
I ι Iota i Φ φ ϕ Phi ph
K κ Kappa k X χ Chi ch
Λ λ Lambda l Ψ φ Psi ps
M µ Mu m Ω ω Omega o
Contents
1 Non-dimensionalisation and approximation 5
1.1 Non-dimensionalisation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 5
1.1.1 The wave equation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 5
1.1.2 Scaling and rescaling . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 7
1.2 Example: projectile problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 7
1.2.1 Scaling the problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 8
2 Introduction to Asymptotics 9
2.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 9
2.2 Useful results . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 9
2.2.1 Theorem 1.1 (Taylor series with remainder) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 9
2.2.2 Theorem 1.2 (l’Hospital’s rule) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 10
2.2.3 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11
2.3 Order Symbols . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11
2.3.1 Definitions: . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11
2.3.2 Examples (for ε ↓ 0) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 12
2.4 Asymptotic Approximations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 12
1
2.4.1 Definition . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 12
2.5 Asymptotic Expansions (Series) and Gauge Functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 13
2.5.1 Definitions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 14
2.5.2 Accuracy and convergence of an asymptotic series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 17
2.6 Stokes phenomenon in the complex plane . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 18
2.6.1 Manipulating Asymptotic Expansions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
3 Algebraic and Transcendental Equations 20

3.1 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 20
3.2 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 21
3.3 Example (A singular perturbation) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 22
3.3.1 Rescaling the equation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 22
3.4 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 23
3.5 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 24
3.6 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 24
3.7 Iterative methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 27
4 Regular Perturbation Solution of Differential Equations 28

4.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 28
4.1.1 Example: Projectiles . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 28
4.1.2 Example: A PDE problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 29
4.1.3 Exercise . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 31
4.1.4 Perturbed boundary conditions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 31
4.2 Uniformity of approximations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 32
4.2.1 Definition of uniform asymptotic approximation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 32
4.2.2 Definition of uniform asymptotic expansion or series . . . . . . . . . . . . . . . . . . . . . . . . . . . . 32
4.3 Small and large time solutions; artificial parameters . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 35
4.4 Maple . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 38
4.5 The chain rule for partial derivatives . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 39
4.6 Solutions to some elementary ODEs. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 40
5 Matched Asymptotic Expansions 42

5.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
5.2 Rescaling differential equations: general considerations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 42
5.3 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 43
5.3.1 Step 1: Outer Solution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 43
5.3.2 Step 2: Rescaling the problem: the boundary layer . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 44
5.3.3 Step 4: Composite expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 48
5.3.4 Step 5: Second term . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 49
5.3.5 Higher order terms . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 49
5.3.6 Discussion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 51
5.4 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 51
5.4.1 Approximate solution using the method of matched asymptotic expansions . . . . . . . . . . . . . . . 52
5.4.2 Matching . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 53
5.5 Multiple boundary layers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 54
5.5.1 Step 1: Outer Expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 54
5.5.2 Steps 2 and 3: Boundary Layers and Matching . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 54
5.5.3 Step 4: Composite Expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 55
5.5.4 Comments . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 56
5.6 Interior Layers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 57
5.6.1 Step 1: Outer Expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 57
5.6.2 Steps 2 and 3: Interior Layer and Matching . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 57
5.7 Corner Layers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 61
5.7.1 Step 1: Outer Expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 61
5.7.2 Step 2: Corner Layer . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 61
5.7.3 Step 3: van Dyke matching . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 62
5.8 Partial Differential Equations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 64
5.8.1 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 64
5.8.2 A special case . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 66
5.9 Extra problems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 73
2
5.9.1 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 73
5.9.2 Example: An elliptic PDE . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 73
5.9.3 Leibniz’ rule for differentiating under the integral sign . . . . . . . . . . . . . . . . . . . . . . . . . . . 78
5.9.4 Travelling periodic waves and exponential solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 79
6 Multiple Scales 79
6.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 79
6.2 A note on resonance . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 79
6.3 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 81
6.3.1 Regular perturbation approach . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 81
6.3.2 Multiple scales approach . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 81
6.4 Derivation of oscillator equation with weak damping . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 83
6.5 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 85
6.5.1 Regular Expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 85
6.5.2 Multiple-Scale Expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 86
6.5.3 Three Time Scales . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 88
6.5.4 Discussion and Observations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 89
6.6 Slowly Varying Coefficients . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 90
6.7 Boundary Layers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 92
6.8 Partial Differential Equations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 93
6.8.1 A linear wave problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 95
6.9 Other multiple scale examples . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 97
6.9.1 Forced Motion Near Resonance . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 97
6.9.2 Nonlinear Waves . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 100
7 Method of strained co-ordinates 106

7.1 A model problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 107
7.1.1 Solution by strained co-ordinates . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 108
7.1.2 Solution by method of matched asymptotic expansions . . . . . . . . . . . . . . . . . . . . . . . . . . . 109
8 The WKB Method 110

8.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 110
8.2 The method . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 110
8.2.1 Examples . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 111
8.2.2 General Discussion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 114
8.3 Turning Points (WKB1-boundary layer-WKB2) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 115
8.3.1 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 115
8.3.2 Examples . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 118
8.4 PDEs modelling wave propagation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 118
8.5 WKB appendix . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 120
8.5.1 Other equivalent WKB substitutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 120
8.5.2 Second Term of a WKB Expansion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 121
8.5.3 Turning point problems: the general case . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 122
8.5.4 Solution in Transition Layer . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 123
8.5.5 Matching . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 123
8.5.6 Summary . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 125
8.5.7 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 126
8.5.8 Energy Methods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 127
8.5.9 Slender Body Approximations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 128
9 Asymptotic expansion of integrals 132

9.1 Preamble . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 132
9.2 A useful result . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 132
9.3 Expansion of integrands . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 133
9.3.1 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 133
9.3.2 Splitting the range of integration: leading order estimates . . . . . . . . . . . . . . . . . . . . . . . . . 133
9.4 Integration by Parts . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 134
9.4.1 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 134
9.4.2 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 135
9.4.3 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 137
9.5 Laplace integrals and Laplace’s Method . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 137
9.5.1 The method in more detail . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 139
3
9.6 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 139
9.6.1 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 140
9.6.2 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 141
9.6.3 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 141
9.6.4 Example . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 142
9.7 Discussion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 142
9.8 Method of Steepest Descents (Saddle point method) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 144
9.8.1 Example Hankel’s integral . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 151
10 Appendix 154
10.1 Gamma Function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 154
10.2 Error function and complementary error function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 154
10.3 Airy Functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 154
10.3.1 Airy Differential equation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 154
10.3.2 General solution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 154
10.3.3 Particular values . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 156
10.3.4 Asymptotic expressions for large . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 156
10.3.5 Useful Taylor (Maclaurin) series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 156
10.4 Binomial series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 157
10.5 Useful Inequalities . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 157
10.5.1 Cauchy Schwartz . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 157
10.5.2 Fredholm alternative . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 158
4
1 Non-dimensionalisation and approximation
Once we have a mathematical model, we have to try and solve it. There are two kinds of solutions: exact, analytical solutions,
and approximate solutions. Exact solutions are explicit formulae; for example we can exactly solve quadratic equations, and
certain differential equations, such as those describing simple harmonic motion. We also consider that solutions such as
Taylor series constitute analytic solutions: they can be computed to arbitrary accuracy. The same applies to quadratures,
such as the solution of
du
= f (u), u(0) = u0 , (1)
dt
which has an implicitly defined solution
Z u
dξ
= t. (2)
0 f (ξ)
There are two methods for obtaining precise approximations: numerical methods and asymptotic or perturbation methods.
Approximate solutions are those where one solves an approximate equation or an approximating sequence of equations.
Asymptotic and perturbation methods work when some parameter is small or large. Numerical methods work best when all
parameters are order one. In practice the two methods complement one another.
In section 4.6 we review solution techniques for some simple ODEs.
1.1 Non-dimensionalisation
In order to approximate a solution, we need to be able to neglect terms which are small. This raises a concept of fundamental
importance, which is that “small” and “large” are adjectives which can only apply quantitatively when a comparison is made
between quantities of the same dimension. An equivalent assertion is that we can only make approximations based on the
small size of parameters if these parameters are dimensionless. It makes no intrinsic sense to say a quantity is small if it still
has dimensions. A speed of one cm s−1 is small if you are a rocket, but large if you are a giant African land snail. An ice
sheet is thick if you are a human being, but thin if you are a planet. So we always associate large or small quantities with
dimensionless numbers, that is, parameters which measure quantities with respect to some reference value with the same
units. The process of isolating these parameters is called non-dimensionalisation.
1.1.1 The wave equation

Putting a mathematical model into non-dimensional form is fundamental. Although technically trivial, there is a certain art
to the process of non-dimensionalisation, and the associated concept of scaling, and the only real way to learn how to do it
is by example. Let us begin with a simple model, the wave equation, which one learns how to derive in first year applied
mathematics courses. We suppose a strong elastic string (fig. 1), for example a guitar string, is wound between two points
a distance l apart. If the tension in the string is T and its density (mass per unit length in this instance) is ρ, then an
application of Newton’s second law to an infinitesimal segment of the string leads to the equation
∂2u ∂2u
ρ 2
= T 2 , u = 0 on x = 0, l, (3)
∂t ∂x
where x is the distance along the string, and u is its transverse displacement.
The main assumption that is usually stated in deriving this equation is that the displacement u is small. However, there
are at least two other implicit assumptions which are made. One is that gravity is unimportant; the other is that there is no
longitudinal displacement.
For a guitar string, these seem to be reasonable assumptions, but why? We expect the effect of gravity to be deviation
of the displacement from the vertical, and this is evidently valid for the guitar string. It is not valid for the hanging chain,
or for the wire between telegraph poles. Why? One would say, for the chain, the density is too large; for the telegraph wire,
the distance l is too large; while the guitar string is straight because it is tight: T is large. These facts suggest that the
“size” of the gravitational acceleration g may in fact be measured by the dimensionless parameter ρgl/T , which appears to
be the only independent dimensionless parameter which can be formed from those in the model if we include gravity.
How can this suspicion be confirmed? From first principles, we derive the wave equation, including gravity, in the form
∂ 2u ∂ 2u
ρ = T − ρg, u = 0 on x = 0, l; u(x, 0) = f (x); ut (x, 0) = h(x). (4)
∂t2 ∂x2
where ρg is the weight force per unit length. Next we write the model in dimensionless form. We do this by non-
dimensionalising the variables so that they are numerically of order one (written O(1)). To do this, we list the variables and
parameters (these are the given quantities ) of the problem and their associated dimensions:
5
u
Figure 1: An oscillating elastic string
variables units
dependent variable u m
independent variable x m
independent variable t s
parameters
density (mass per unit length) ρ Kg m−1
−2
tension T N ≡ Kg m s
gravity g m s −2
length l m
max initial displacement u0 = max(f (x)) m
We include the maximum initial displacement u0 in this list so that we can non-dimensionalise the initial condition. (We could
also include a measure of the initial velocity h(x) (which has units ms−1 ) among these parameters). To non-dimensionalise,
we need to choose parameters, or combinations of parameters which have the same units as the variables. There are
many ways to do this. In the present case we make choices in order to balance terms in the equation. Specifically, in order
to balance the first two terms in the PDE we write
x = lx∗ , u = u0 u∗ , t = (l/c)t∗ , (5)
where
s
T
c≡ (6)
ρ
is the wave speed, and u0 is a measure of the displacement: for example, here it is the maximum initial displacement.
q Note
2
that the time t could only be nondimensionalised by using a combination of the parameters. (Check that l/c ≡ ρlT does
indeed have the dimensions of time). Thus we transform u(x, t) → u∗ (x∗ , t∗ ). To transform derivatives we use the chain rule,
remembering that the parameters are all constants. Thus, for example,
∂u ∂(u0 u∗ ) u0 ∂u∗ ∂2u u ∂ 2 u∗ ∂u ∂(u0 u∗ ) u0 c ∂u∗

≡ ∗
≡ ; ≡ 20 ∗2 ; ≡ l
≡ . (7)
∂x ∂(lx ) l ∂x∗ ∂x2 l ∂x ∂t ∗
∂(( c )t ) l ∂t∗
Note that all the dimensional information is contained in the premultipliers, the starred terms are all dimensionless. Thus,
u0
one should always check to ensure that ∂u
∂x has the same dimensions as l etc. Proceeding in this fashion, the dimensionless
model is then obtained in the form
∂ 2 u∗ ∂ 2 u∗
∗2
= − β, (8)
∂t ∂x∗2
where we have a single dimensionless parameter
ρgl2
β= , (9)
T u0
6
with u∗ = O(1) initially. The boundary and initial conditions must also be transformed and this gives:
f (x) ∂u∗ (x∗ , 0) g(x)l

u∗ (x∗ = 0, t∗ ) = 0; u∗ (x∗ = 1, t∗ ) = 0; u∗ (x∗ , 0) =; ∗
= . (10)
u0 ∂t u0 c
All of the terms in the equations and in the initial and boundary conditions are dimensionless, and all the coefficients which
appear (such as β) are dimensionless.
It is conventional at this dimensionless stage to dispense with the asterisks in writing the variables, and this we now do.
The process of choosing particular scales for the variables, or scaling, is motivated by the following ideas. Firstly, there is
a natural length scale to the problem, l, which is the dimension of the geometric domain on which the problem is to be
resolved. Further, there is a natural length scale for the displacement, u0 , which is present in the initial conditions. Next,
of the three terms in (4), we anticipate (at least for the guitar string) that the gravity term will be “small”. It follows
that the other two should be the same size, and we choose the time scale so that these two terms “balance”, that is their
dimensional scales (here T u20 /l2 ) are the same. When the model is written in the dimensionless form (8), we then have equal
dimensionless coefficients multiplying these two terms: here they are both equal to one.
The final, essential idea is that, in general, a dimensionless function u(x, t) which varies by O(1) over an x range of O(1)
will have derivatives of O(1). This is true, for example, for sin x, e−x , and x2 : it is not true for the function e−10x , which
varies rapidly over a distance of order x ∼ 0.1 near x = 0. With this assumption, the derivative terms ∂ 2 u/∂t2 and ∂ 2 u/∂x2
are O(1), and it follows that the relative size of the gravity term is given by β. Thus gravity is negligible if β ≪ 1, and indeed
this means large tension, small density or short length, as we surmised.
1.1.2 Scaling and rescaling

The key idea behind a good nondimensionalisation is that it should make explicit the balances evident in a particular
problem i.e., it should allow one to identify (via the occurrence of dimensionless parameters whose sizes are known) the
more important terms before one actually solves the problem. The problem is then said to be well scaled. There will also
be situations where one set of scales cannot reflect the overall behaviour and two or more sets are necessary; in such cases
we have a non-regular or singular perturbation and we will speak of rescaling the problem to demonstrate the different
balances. In these notes, the problems we deal with will all be assumed to be in a well scaled form i.e., we assume a
suitable nondimensionalisation has already been carried out reflecting at least one balance of significance to
the problem.
Note that from an exact mathematical point of view, two different scalings of a particular problem are equivalent. If we
can solve a problem exactly, it does not matter which scaling we use. However, from an asymptotic point of view, if we
wish to exploit the existence of small parameters to find approximate solutions, then different scalings will represent quite
different asymptotic balances.
1.2 Example: projectile problem

A simple example comes from the motion of an object projected radially upward from the surface of the Earth. Letting x(t)
denoted the height of the object, measured from the surface, from Newton’s second law we obtain the following equation of
motion:
d2 x Mm
m = −G (11)
dt2 (x + R)2
where m, M are the mass of the object and the earth and G is the proportionality constant in the universal gravitational law.
When x = 0, at the earth’s surface, the gravitational force must equal mg and so GM/R2 = g where g is the acceleration
due to gravity on the surface of the earth. Thus we have
d2 x gR2
2
=− , for 0 < t (12)
dt (x + R)2
where R is the radius of the Earth and g is the gratitational constant. We will assume the object starts from the surface with
a given upward velocity, and so x(0) = 0 and x′ (0) = v0 where v0 is positive. The nonlinear nature of the above ordinary
differential equation makes finding a closed-form solution difficult, and it is thus natural to try to find some way to simplify
the equation. For example, if the object does not get far from the surface, then one might try to argue that x is small
compared to R and the denominator in (12) can be simplified to R2 . This is the type of argument often made in introductory
′′ ′
physics and engineering texts. In this case x ≈ x0 , where x0 = −g for x0 (0) = 0 and x0 (0) = v0 .
The solution in fig. 2 comes from the linearization of (12) and corresponds to the motion in a uniform gravitational field.
This solution is
1
x0 (t) = − gt2 + v0 t. (13)
2
One finds in this case that the object reaches a maximum height of v02 /2g and comes back to Earth when t = 2v0 /g (see Fig
2). The difficulty with this reduction is that it is unclear how to determine a correction to the approximate solution in (13).
7
x0
v02/2g
v0/g t
Figure 2: Schematic of the solution x0 (t) given in (13). This solution comes from the linearization of (12) and corresponds
to motion in a uniform gravitational field.
This is worth knowing since we would then be able to get a measure of the error made in using (13) as an approximation
and it would also be possible to see just how the nonlinear nature of the original problem affects the motion of the object.
1.2.1 Scaling the problem

To make the reduction process more systematic, it is first necessary to scale the variables. To do this, let t∗ = t/tc and
x∗ (t∗ ) = x(t)/xc , where tc is a constant characteristic time for the problem and xc is a constant characteristic value for
the solution (both of which we choose from a consideration of the (constant) parameters in the problem). We have a lot of
freedom in picking these constants, but they should be representative of the situation under consideration.
variables units
dependent variable x m
parameters
gravity g ms−2
radius of earth R m
initial velocity v0 ms−1
From the parameters, we need to choose combinations which have the same units as the variables and which lead to
a well scaled problem i.e., which give a real representation of the way the equation behaves from an asymptotic point of
view. Note, for example, that choosing xc = R would be a poor choice as typically one will not be throwing a ball a distance
equal to the radius of the earth R. Based on what is shown in fig. 2, we take the time scale to be tc ≡ v0 /g and the length
scale to be xc ≡ v02 /g. Doing this, (using the chain rule for the derivatives)
v2
dx d( g0 x∗ ) dx∗ dx2 d2 x∗
≡ v0 ∗ ≡ v0 ; ≡ g ∗2 ; (14)
dt d(( g )t ) dt∗ dt 2 dt
the problem transforms into the following:
d2 x∗ 1
=− , 0 < t∗ . (15)
dt∗2 (1 + εx∗ )2
The boundary conditions must also be transformed. This gives:
dx∗
x∗ (0) = 0; (0) = 1. (16)
dt∗
In (15), the parameter ε ≡ v02 /Rg is dimensionless and its value is important because it gives us a measure of how high the
projectile gets in comparison to the radius of the Earth. In terms of the function x0 (t), it can be seen from fig. 2 that ε/2
is the ration of the maximum height of the projectile to the radius of the Earth. Assuming R = 6000 km, g = 10 ms−2 then
v02 −2 2
ε ≈ 6.10 7m s . It would therefore appear that if ε ≪ 1 i.e. if v0 is much less than 8 × 103 ms−1 , then (13) is a reasonably
good approximation to the solution. We can verify this assertion by reconstructing the first-order approximation in (13).
This can be done by assuming that the dependence of the solution on ε can be determined using a Taylor series expansion
about ε = 0. In other words, for small ε, it is assumed that
8
x∗ ∼ x∗0 (t∗ ) + εx∗1 (t∗ ) + · · · .
The first term in this expansion is the scaled version of x0 , (this will be shown later). What is important is that with this
approach it is possible to estimate how well x∗0 approximates the solution of (15,16) by finding x∗1 . The method of deriving the
zero-order approximation (x∗0 ) and its correction (x∗1 ) is not difficult, but we first need to put the definition of an asymptotic
approximation on a firm foundation. Those interested in investigating the ideas underlying the nondimensionalization of a
physical problem, as well as some of the theory underlying dimensional analysis, should consult Lin and Segel (1989).
Choosing scales is about balances In fact, the whole scaling process just amounts to choosing the correct
balances in the problem. In this problem we have to choose two scales: one for x∗ and one for t∗ . This suggests that we
must find two fundamental balances associated with this problem. Specifically if we reconsider (12), we note that because
of the “largeness” of R, we expect that x∗ ≪ R so that x∗ + R ≈ R and the RHS can be approximated by −gR2 /R2 = −g.
Hence the approximate balance should be
d2 x∗
≈ −g
dt∗2
If we scale as follows:
x∗ = xL; t∗ = tT,
where L, T have not yet been chosen then the approximate balance becomes
L d2 x
≈ −g
T 2 dt2
and this suggests that we must choose:
L
= g =⇒ L = gT 2 . (17)
T2
If we now examine the initial condition in dimensionless form:
L dx
(t = 0) = V
T dt
and this suggests that we should choose
L L
= V =⇒ T = (18)
T V
Putting (17) and (18) together, we obtain:
L2 V2 V
L=g 2
=⇒ L = =⇒ T = . (19)
V g g
2 Introduction to Asymptotics
2.1 Introduction
In these notes, we are interested in using what are known as asymptotic expansions to find approximate solutions of differential
equations. Usually our efforts will be directed toward constructing the solution of a problem with only occasional regard for
the physical situation it represents.
2.2 Useful results

As in the preceding example, we will typically end up expanding functions in powers of ε. Given a function f (ε), one of the
most important tools for doing this is Taylor’s theorem. This is a well-known result, but for completeness it is stated below.
2.2.1 Theorem 1.1 (Taylor series with remainder)

Given a function f (ε), suppose its (n + 1)st derivative f (n+1) is continuous for εa < ε < εb . In this case, if ε0 and ε are points
in the interval (εa , εb ), then
1
f (ε) = f (ε0 ) + (ε − ε0 )f ′ (ε0 ) + · · · + (ε − ε0 )n f (n) (ε0 ) + Rn+1 , (20)
n!
where
1
Rn+1 = (ε − ε0 )n+1 f (n+1) (ξ), (21)
(n + 1)!
9
and ξ is a point between ε0 and ε.
The student may be more familiar with Taylor series in the form:
1 n (n)
f (x0 + h) = f (x0 ) + hf ′ (x0 ) + · · · + h f (x0 ) + Rn+1 ,
n!
which is equivalent with h ≡ ε − ε0 , x ≡ ε0 . If one wishes to write a Taylor series for f (x) about the point x0 = π/4 then
we write
1
x = x0 + h = π/4 + h =⇒ h = x − π/4 =⇒ f (x) = f (π/4) + (x − π/4)f ′ (x − π/4) + (x − π/4)2 f ′′ (π/4) . . .
2
The special case where x0 = 0 is sometimes referred to as the MacLaurin series.
This result is useful because if the first n + 1 terms from the Taylor series are used as a approximation of f (ε), then it is
possible to estimate the error using (21).
The Taylor series probably used most often in these notes is the binomial expansion, which states
1
(x + y)α = xα + αyxα−1 + α(α − 1)xα−2 y 2 + ....
2!
This holds for any real α and y 2 < x2 . Other expansions that we will often have use for include the Maclaurin series
(ε0 = 0) of the exponential and trigonometric functions. (See end of notes for some useful results).
Another useful result is l’Hospital’s rule, which concerns the value of the limit of the ratio of two functions in indeterminate
situations.
2.2.2 Theorem 1.2 (l’Hospital’s rule)

Suppose f (ε) and φ(ε) are differentiable on the interval (ε0 , εb ) and φ′ (ε) 6= 0 in this interval. Also suppose
f ′ (ε)
lim = A,
ε↓ε0 φ′ (ε)
where −∞ ≤ A ≤ ∞. In this case,
f (ε)
lim = A,
ε↓ε0 φ(ε)
if either one of the following conditions holds:
1. f → 0 and φ → 0 as ε ↓ ε0 ,
2. φ → ∞ as ε ↓ ε0 . N.B.
l’Hospital’s rule need only be used where the limit when evaluated directly is indeterminate (when a 0/0 or ∞/∞
situation arises. Situations where a limit evaluates to 1/0 = ∞, 1/∞ = 0, 0/1 = 0 are not considered indeterminate).
For example consider the case where f = cos ε, φ = ε as ε → 0. Then
f cos ε limε→0 cos ε 1
lim = lim = → →∞ (22)
ε→0 g ε→0 ε limε→0 ε limε→0 ε
(− sin ε)
But if we (incorrectly) use l’Hospital’s rule we would get the incorrect result limε→0 1 = 0. For the case f =
ε2 , φ = ε we can use l’Hospitals rule, but it is simpler to simplify algebraically:
f ε2
lim = lim = lim ε = 0. (23)
ε→0 g ε→0 ε ε→0
Recall also that if a limit gives rise to a C/∞ situation, where C is a finite constant, the limit of course corresponds to
0 and there is no need to use l’Hospital’s rule (though it is not incorrect to do so). For example if f = 1, g = 1/ε then
f 1 1 limε→0 ε
lim = lim (= )= = 0.
ε→0 g ε→0 1/ε ∞ 1
The proofs of these theorems, and some of their consequences, can be found in Rudin (1976).
10
2.2.3 Example
Use l’Hospital’s rule to prove that for α > 0, xα ln x → 0 as x → 0.
0 ∞ ln x ∞
As x → 0, xα → 0, ln x → −∞. To use l’Hospital’s rule we need a 0 or ∞ situation. Thus write xα ln x = x−α which is a ∞
case. Using l’Hospital we thus have:
ln x x−1 xα
lim = lim = lim (− ) = 0.
x→0 x−α x→0 −αx−α−1 x→0 α
2.3 Order Symbols

To define an asymptotic approximation, we first need to introduce order, or Landau, symbols. The reason for this is that we
will be interested in how functions behave as a parameter, typically ε, becomes small. For example, the function φ(ε) = ε
does not converge to zero as fast as f (ε) = ε2 when ε → 0, (see fig. 4) and we need a notation to denote this fact. To
understand what follows, remember that in asymptotics we are interested in mathematical behaviour for some particular
limiting behaviour, most often as some parameter tends to 0 or ∞.
2.3.1 Definitions:
If the limit f (ε)/φ(ε) exists (including the case where f /φ → ∞) then we say that:
1. f = O(φ) as ε ↓ ε0 if the limit

f (ε)
lim (24)
ε↓ε0 φ(ε)
is finite. We then say that f is “big Oh” of φ as ε ↓ ε0 or that f is asymptotically bounded by g in magnitude as
ε → ε0 . We usually think of this as meaning that f behaves like a constant times φ for ε sufficiently close to ε0 .
2. f = o(φ) as ε ↓ ε0 if
f (ε)
lim =0 (25)
ε↓ε0 φ(ε)
and we say that f is “little oh” of φ as ε ↓ ε0 or that f is asymptotically smaller than g as ε → ε0 . We think of this as
meaning that f is much smaller than φ for ε sufficiently close to ε0 . Sometimes this is written as f ≪ φ.
Note that f = o(φ) =⇒ f = O(φ) but this is only a one way implication.
3. f ∼ φ as ε ↓ ε0 if f − φ = o(φ) or equivalently
f (ε)
lim =1 (26)
ε↓ε0 φ(ε)
in which case we say that “f is asymptotic to φ” or “f twiddles φ”. Note that this is a special case of f = O(φ).
4. f = Os (φ) as ε ↓ ε0 ⇐⇒ f = O(φ) but f 6= o(φ). This notation is not very commonly used but can be very useful.
In this case we say that “f is strictly of order φ”.
Aside The above definitions are most often used in practice and we will be satisfied with them. More precise definitions
for the case are included here for completeness:
1. f = O(φ) as ε ↓ ε0 means that there are constants k0 and ε1 (independent of ε) such that
|f (ε)| ≤ k0 |φ(ε)| for ε0 < ε < ε1 .
2. f = o(φ) as ε ↓ ε0 means that for every positive δ there is an ε2 (independent of ε) such that
|f (ε)| ≤ δ|φ(ε)| for ε0 < ε < ε2 .
This is also commonly written as f ≪ φ.

f −φ
3. f ∼ φ as ε ↓ ε0 means that f − φ = o(φ) as ε ↓ ε0 . Note that if the limits exist this is equivalent to saying that φ →0
f
i.e. φ → 1.
4. If g 6= 0, then f = O(g) as ε → ε0 if and only if f /g is bounded in a (punctured) neighborhood of ε0 , and f = o(g) if

and only if f /g → 0 as x → 0.
These definitions may seem cumbersome, but they usually are not hard to apply. The proofs of these statements follow
directly from the definition of a limit and are left to the reader.
11
2.3.2 Examples (for ε ↓ 0)
1. Suppose f = ε2 . Also let φ1 = ε and φ2 = −3ε2 + 5ε6 . In this case,
f
lim = 0 ⇒ f = o(φ1 )
ε↓0 φ1
and
f
lim = −1/3 ⇒ f = O(φ2 ).
ε↓0 φ2
2. Show that ∼ defines an equivalence relationship i.e., f ∼ f , f ∼ g ⇒ g ∼ f and f ∼ g, g ∼ h ⇒ f ∼ h.
3. If f (ε) = sin(ε) then, using Taylor’s theorem, f = ε − 12 ε2 sin(ξ).Thus, limε↓0 (f /ε) = 1, and from this it follows that
f = O(ε) and f ∼ ε.
4. If f = e−1/ε then in fact f = o(εα ) as ε → 0 for all values of α. To see this, it is simplest to put ε = 1/δ and consider
the limit δ → ∞. For the case α = 1 we then have
dδ
e−1/ε δ 1
lim ≡ lim δe−δ = lim δ = lim dδ
deδ
= lim =0 (27)
ε→0 ε δ→∞ δ→∞ e δ→∞
dδ
δ→∞ eδ
using L’Hospital’s rule for the last step. A similar result can be obtained for any positive α.
We say in this case that f is transcendentally or exponentially small with respect to the power functions εα (See
also fig. 4). The concept of being exponentially small is, of course, a relative one. exp(−3/ε) is not for example
exponentially small wrt exp(−1/ε), ˙ exp(−2/ε), exp(−3/ε)... Similarly if one happens to be using the set of well-ordered
gauge functions {ε ln ε, ε (ln ε) , ε (ln ε)3 ...} then εα is exponentially small wrt this set because εα ≪ εα (ln ε)α ∀α ∈ N.
2 2 3
5. If f = ε sin(1 + 1/ε) and φ = ε. In this case we cannot use the limiting definition (the limit does not exist). Instead we
use the fact that |f /φ| ≤ 1 for 0 < ε. Hence, f = O(φ).
Some of the properties of the order symbols are examined in the exercises. Three that are worth mentioning are the
following (the symbol ⇐⇒ stands for the statement “if and only if”):
6. f = O(1) as ε ↓ ε0 ⇐⇒ f is bounded as ε ↓ ε0 .
7. f = o(1) as ε ↓ ε0 ⇐⇒ f → 0 as ε ↓ ε0 .
8. f = o(φ) as ε ↓ ε0 =⇒ f = O(φ) as ε ↓ ε0 (but not vice versa).
1
9. Prove (exercise) that − ln ε ≡ ln (1/ε) ≪ εα for any α > 0. (“Powers beat logs”).
α
Another way of expressing this is ε ln (1/ε) → 0 as ε → 0.
The proofs of these statements are straightforward.
Two symbols we will also use occasionally are ≪ and ≈ . When saying that f (ε) ≪ φ(ε) it is meant that f = o(φ), and
the statement that ε ≪ 1, or that ε is small, means ε ↓ 0. We will also use x = O(1) and x ∼ 1 occasionally in a loose sense
to mean that x is typically of about magnitude 1 i.e. x is well scaled. The symbol ≈ does not have a precise definition and
is used simply to designate an approximate numerical value. An example of this is the statement that π ≈ 3.14.
2.4 Asymptotic Approximations

Our objective is to construct approximations to the solutions of equations. It is therefore important that we state exactly
what we mean by an approximation. To introduce this idea, suppose we are interested in finding an approximation of
f (ε) = ε2 + ε5 for ε close to zero. Because ε5 ≪ ε2 , a reasonable approximation is f (ε) ≈ ε2 . On the other hand, a lousy
approximation is f (ε) ≈ 32 ε2 . This is lousy even though the error f (ε) − 23 ε2 goes to zero as ε ↓ 0. The problem with
this “lousy approximation” is that the error is of the same order as the function we are using to approximate f (ε). This
observation gives rise to the following definition.
2.4.1 Definition
Given f (ε) and φ(ε), we say that φ(ε) is an asymptotic approximation to f (ε) as ε ↓ ε0 whenever f = φ+o(φ) for ε sufficiently
close to ε0 (or equivalently lim φf = 1). In this case, we write f ∼ φ as ε ↓ ε0 .
As demonstrated in the above example, the idea underlying this definition is that φ(ε) serves as an approximation to
f (ε), for ε close to ε0 , when the error is of higher order i.e., much smaller, than the approximating function. In the case
where φ(ε) is not zero near ε0 , we can make use of (25). In particular, we have that f ∼ φ as ε ↓ ε0 if
f (ε)
lim = 1. (28)
ε↓ε0 φ(ε)
12
1.0
0.8
Function
0.6
0.4
0.2 Function
Asymptotic Approximation
0 0.2 0.4 0.6 0.8 1.0
x - axis
Figure 3: Comparison between the function f = x + e−x/ε and its asymptotic approximation f ∼ x. Note that the two
functions are essentially indistinguishable except near x = 0. In the plot ε = 0.01.
Example 1 Suppose f = sin(ε) and ε0 = 0. We can obtain an approximation of f (ε) by expanding about ε = 0 using
Taylor’s theorem. This yields
1 1 5
f = ε − ε3 + ε cos(ζ).
6 120
From this, the following asymptotic approximations are obtained:
i) f ∼ ε, ii) f ∼ ε − 16 ε3 , iii) f ∼ ε + 2ε2 .
In comparing these asymptotic approximations with the Taylor series expansion, it would appear that, for small ε, (ii) is
the most accurate and (iii) is the least accurate. However, our definition of an asymptotic approximation says little about
comparative accuracy. We correct this weakness later by introducing asymptotic expansions.
Example 2 Suppose f = x + e−x/ε where 0 < x < 1 is fixed. In this case f ∼ x for small ε. However, it is natural to ask
how well this does in approximating the function for 0 < x < 1. If we plot both together, we obtain the curves shown in
Fig. 3. It’s apparent that the approximation is quite good away from x = 0. It is also clear that we do not do so well near
x = 0. This is true no matter what value of ε we choose since f (0) = 1. It should be remembered that with the asymptotic
approximation, given a value of x, the approximation is a good one if ε is close to ε0 = 0. What we are seeing in this example
is that exactly how small ε can be depends on x (the closer we are to x = 0 the smaller ε has to be). In later sections, we
will refer to this situation as a case of when the approximation is not uniformly valid on the interval 0 < x < 1.
Example 3 Consider the function f = sin(πx) + ε3 for 0 ≤ x ≤ 12 . For small ε it might seem reasonable to expect
that f ∼ sin(πx). For this to hold, it is required that f − sin(πx) = o(sin(πx)) as ε ↓ 0. If x 6= 0, this is true since
limε↓0 (ε3 / sin(πx)) = 0. However, at x = 0 this requirement does not hold since sin(0) = 0. Therefore, sin(πx) is not an
asymptotic approximation of f over the entire interval 0 ≤ x ≤ 12 . This problem of using an approximating function is one
that we will come across on numerous occasions. Usually, it is a complication we will not worry a great deal about since the
correction is relatively small. For instance, in this example the correction is O(ε3 ), while everywhere else the function is O(1).
This is not true for the correction that is needed to fix the approximation in the previous example. There the correction at,
or very near, x = 0 is O(1), and this is the same order as the value of the function through the rest of the interval. We are
therefore not able to ignore the problem at x = 0, and we will deal with this later.
2.5 Asymptotic Expansions (Series) and Gauge Functions

Two observations that come from the preceding examples are that an asymptotic approximation is not unique, and it also
does not say much about the accuracy of the approximation. To address these shortcomings, we need to introduce more
structure into the formulation. In the earlier examples, the sequence 1, ε, ε2 , ε3 , ... was used in the expansion of the function
but other, usually more interesting, sequences also arise. In preparation for this, we have the next definitions, which are due
to Poincaré (1886).
13
2.5.1 Definitions
1. The functions φ1 , φ2 , ...form an asymptotic sequence, or are well-ordered, as ε ↓ ε0 if and only if φn = o(φm ) as ε ↓ ε0
for all m and n that satisfy m < n.
2. If φ1 , φ2 , ... is an asymptotic sequence, then f (ε) has an asymptotic expansion to n terms, with respect to this sequence,
if and only if
m
X
f= ak φk (ε) + o(φm ) for m = 1, ..., n as ε ↓ ε0 , (29)
k=1
where the ak s are independent of ε. This means that (for ε sufficiently close to ε0 )
f = a1 φ1 (ε) + o(φ1 ) and f = a1 φ1 (ε) + a2 φ2 (ε) + o(φ2 ) etc.
In terms of limits this means that

f − a1 φ1 (ε) f − a1 φ1 (ε) − a2 φ2 (ε)
lim = 0 and lim = 0 etc.
ε→ε0 φ1 ε→ε0 φ2
The φk s are called the scale, or gauge, or basis functions (see fig. 4).
To make use of this definition, we need to have some idea of what scale functions are available. We will run across a wide
variety in these notes, but a couple of our favourites will turn out to be the following:
1. φ1 = εα , φ2 = εβ , φ3 = εγ , · · · where α < β < γ < · · · (for ε ↓ 0), and

2. φ1 = 1, φ2 = e−1/ε , φ3 = e−2/ε , . . . (for ε ↓ 0).
The first of these is simply a generalization of the power series functions, and the second sequence is useful when we
have to describe exponentially small functions. The verification that these do indeed form asymptotic sequences is left to
the reader.
ε3
Note that if f = sin ε, φ = ε + 2ε3 , then f ∼ φ but φ is not a two term asymptotic expansion for f . But φ1 = ε − 3! is
a two term asymptotic expansion for f with error o(ε3 ).
Asymptotic expansion of a known function f (ε)

Now comes the harder question. Given a function f (ε), how do we find an asymptotic expansion of it? The most commonly
used methods include employing one of the following
1. Taylor’s theorem,
2. l’Hospital’s rule,
3. an educated guess.
The latter usually relies on an intuitive understanding of the problem and, many times, on luck. The other two methods are
more routine and are illustrated in the examples ahead. The hardest step is usually finding the first term in the asymptotic
series: in particular establishing the relevant set of gauge functions for the problem.
Taylor’s theorem is a particularly useful tool because if the function is smooth enough to let us expand about the point
ε = ε0 , then any one of the resulting Taylor polynomials can be used as an asymptotic expansion (for ε ↓ ε0 ). Moreover,
Taylor’s theorem enables us to analyze the error very easily.
Example using Taylor’s theorem

1. To find a three-term expansion of eε , for ε ≪ 1, we proceed as follows:
1 1
eε = 1 + ε + ε2 + ε3 + · · ·
2 3!
1 2
∼1+ε+ ε
2
ε
To find a three-term expansion of sin(e ), we first use Taylor’s theorem to expand sin(1 + α) about α = 0 to obtain
1
sin(1 + α) = sin(1) + α cos(1) − α2 sin(1) + · · ·
2
14
y=ε1/3
0.9
0.8
y=ε1/2
0.7 y=ε
0.6
0.5 y=ε2
y
0.4
0.3
0.2
0.1 y=exp(-1/ε)
0.0
0.0 0.1 0.2 0.3 0.4 0.5 0.6 0.7
Figure 4: Comparison of sizes of some typical gauge functions as ε → 0. Note that in this limit exp(−1/ε) ≪ εn for any
positive n. It is said to be exponentially (or transcendentally) small.
Therefore
1 1 1 1
sin(eε ) = sin(1 + ε + ε2 + ε3 + . . .) = sin(1) + (ε + ε2 + . . .) cos(1) − (ε + . . .)2 sin(1) + . . . (30)
2 3! 2 2
1 2 3
= sin(1) + ε cos(1) + ε (cos(1) − sin(1)) + O(ε ). (31)
2
Note Here we recognize that the function we are dealing with is a compound function. So we start by expanding the
innermost function eε so we have:
1
eε ∼ 1 + ε + ε2 + . . .
2
From this we conclude that
ε 1 2
sin (e ) ∼ sin 1 + ε + ε + . . .
2
Next, we observe the sine function of the right hand side has the form sin (1 + α) where α = ε + 21 ε2 + . . . is close to zero for
small ε. Hence we can expand the sine function in a Taylor series about 1 so that
2
sin (eε ) ∼ sin (1 + α) ∼ sin 1 + α cos 1 − 12 α2 sin 1 + . . . ∼ sin 1 + cos 1 ε + 21 ε2 + . . . − 1
2 sin 1 ε + 12 ε2 + . . .
∼ sin 1 + ε cos 1 + 12 ε2 (cos 1 − sin 1) + . . .
An alternative approach would be to first expand exactly sin(1 + α) = sin 1 cos α + cos 1 sin α and now to use McLaurin
expansions for sin α and cos α.
A nice aspect of using Taylor’s theorem is that the scale functions do not have to be specified ahead of time. This differs
from the next procedure, which requires the specification of the scale functions before constructing the expansion
Example using l’Hospital’s rule To describe the second procedure for constructing an asymptotic expansion, suppose the
scale functions φ1 , φ2 . . . are given and the expansion of the function is f ∼ a1 φ1 (ε)+a2 φ2 (ε)+. . . . From the above definition,
this means that f = a1 φ1 + o(φ1 ). Assuming that we can divide by φ1 , we get limε↓ε0 (f /φ1 ) = a1 . This gives us the value
of a1 and with this information we can determine a2 by noting f = a1 φ1 + a2 φ2 + o(φ2 ). Thus limε↓ε0 [(f − a1 φ1 )/φ2 ] = a2 .
This idea can be used to calculate the other coefficients of the expansion, and one obtains the following formulas:
f
a1 = lim , (32)
ε↓ε0 φ1
f − a1 φ1
a2 = lim , (33)
ε↓ε0 φ2
15
f − a1 φ1 − a2 φ2
a3 = lim ,··· (34)
ε↓ε0 φ3
This assumes that the scale functions are nonzero for ε near ε0 and that each of the limits exists. If this is the case, then
above formulae for the ak s show that the asymptotic expansion is unique (but only wrt a particular set of gauge functions).
Example using l’Hospital’s rule

Suppose φ1 = 1, φ2 = ε, φ3 = ε2 , . . . and
1
f (ε) = + e−1/ε.
1+ε
From the limit formulas just given, we have that
f
a1 = lim = 1,
ε↓0 1

f −1 −1 1 −1/ε
a2 = lim = lim + e = −1,
ε↓0 ε 1+ε ε
ε↓0

f −1+ε 1 1 −1/ε
a3 = lim = lim + e = 1, . . . .
ε↓0 ε2 ε↓0 1 + ε ε2
Thus, f ∼ 1 − ε + ε2 + . . . . What is interesting here is that the exponential makes absolutely no contribution to this
expansion. This is because e−1/ε = o(εα ) for all values of α; in other words, the exponential decays so quickly to zero that
any function considers it to be zero. (See also fig. 4) In this case, we say that this exponential function is transcendentally or
exponentially small with respect to these scale functions (i.e., exp(−1/ε) is exponentially small in comparison to 1, ε, ε2 . . .)
As seen in the previous example, two functions can have the same asymptotic expansion. In particular, using the power
series functions φ0 = 1, φ1 = ε, φ2 = ε2 , . . . one obtains the same expansion as in the previous example for any of the following
tanh(1/ε)
1. f1 = 1+ε ,
1+e−1/ε
2. f2 = 1+ε ,
1
3. f3 = 1+ε + ε100 sech(−1/ε).
This observation brings up the idea of asymptotic equality, or asymptotic equivalence, with respect to a given sequence
φ1 , φ2 , φ3 , . . . We say that two functions f and g are asymptotically equal, to n terms, if f − g = o(φn ) as ε ↓ ε0 .
Example Find a four term asymptotic expansion for f = 1 + sin ε in terms of the gauge functions {1, ε, ε2 . . .} and
{1, sin ε, ε2 , ε3 . . .}
Solution Note first that both of the sets of gauge functions are well ordered (?) as ε → 0. The second set looks a little
strange, but it is a perfectly valid set of course.
It is easy to see that in terms of the second set {1, sin ε, ε2 , ε3 . . .} that if
1 + sin ε ∼ a0 + a1 sin ε + a2 ε2 + . . .
then we have the exact result a0 = 1, a1 = 1, a2 = a3 = 0 . . . However, in terms of the first set we find that if
1 + sin ε ∼ b0 + b1 ε + b3 ε2 + b4 ε3 + ..
then (using a Taylor series for sin ε about ε = 0) we have
1
b0 = 1, b1 = 1, b2 = 0, b3 = −
3!
so asymptotic expansions are different if we choose different sets of gauge functions. But there is nothing sacrosanct about
the powers of ε as gauge functions. They are commonly used, but sometimes there are other sets which suit a particular
problem better.
16
2.5.2 Accuracy and convergence of an asymptotic series
It is not unusual to expect that to improve the accuracy of an asymptotic expansion it is simply necessary to include more
terms. This is what happens with Taylor series expansions, where in theory one should be able to obtain as accurate a result
as desired by simply adding together enough terms. However, with an asymptotic expansion this is not necessarily true. The
reason is that an asymptotic expansion only makes a statement about the series in the limit of ε ↓ ε0 , whereas increasing the
number of terms is saying something about the series as n → ∞. There are two limiting processes going on and these need
not commute. In fact, an asymptotic expansion need not and very often will not converge! Moreover, even if it converges it
does not have to converge to the function that was expanded! These two observations may seem to be major flaws in our
definition of an asymptotic expansion but they are, in fact, attributes of which we can take advantage.
To fix ideas, suppose a function f (ε) is represented by the first N terms of a series in powers of ε plus a remainder RN (ε)
so that
N
X
f (ε) = an εn + RN (ε) (35)
n=0
where the an are independent of ε. This series converges if and only if

N
!
X
n
lim RN ≡ lim f (ε) − an ε = 0. (36)
N →∞ N →∞
ε fixed ε fixed n=0
Do not confuse this with the necessary condition that the nth term must tend to zero as n → ∞ in order for the series to
converge. Think of this as fixing ε and letting N → ∞. This series is an asymptotic series for ε → 0 if
RN (ε) = o(εN ) (ε → 0), (37)

i.e., the remainder is asymptotically smaller than the last retained term. Think of this as fixing N and letting ε → 0.
Clearly, a convergent series is an asymptotic series; however, an asymptotic series need not converge.
A demonstration that a convergent asymptotic expansion need not converge to the function that was expanded can be found
in example 2.5.1. It is not hard to show that the asymptotic series converges to the function (1 + ε)−1 . This is clearly not
equal to the original function since it is missing the exponential term. In general the divergence of the asymptotic series
representing a function indicates that the function has a singularity at ε = 0.
A well-known example of a divergent asymptotic expansion arises with the Bessel function J0 (z), which is defined via the
convergent series:
∞
X (−1)k z 2k
J0 (z) = . (38)
22k (k!)2
k=0
If we let f (ε) = J0 ( 1ε ), then it can be shown that an asymptotic expansion of f (ε) for small ε is (Abramowitz and Stegun,
1972)
r
2ǫ 1 π 1 π
f∼ α cos − + β sin − , (39)
π ε 4 ε 4
where
1·32 ε2 1·32 ·52 ·72 ε4

α∼1− + + ... (40)
2!82 4!84
and
ε 1·32 ·52 ε3
β∼ − + ... (41)
8 3!83
It is not hard to show that the series expansions in (40) and (41) are divergent for all nonzero values of ε (Exercise). To
see just how well a divergent series like (39) can approximate the function f (ε), the errors in using (38) and (39) are shown
in Fig. 3 for ε = 1/15. What is immediately apparent is that the asymptotic approximation does very well using only one
or two terms, while the series needs to include more than twenty terms to achieve the same error. If a smaller value of ε is
used, the situation is even more pronounced, eventually getting to the point where it is essentially impossible to calculate
the value of the Bessel function using the series representation since so many terms are needed.
These observations concerning convergence should always be kept in mind when using asymptotic expansions. Also, it
should be pointed out that asymptotic approximations are most valuable when they can give insights into the structure of
the solution, or the physics of the problem under study. If decimal-place accuracy is desired, then numerical methods, or
perhaps numerical methods combined with an asymptotic solution, are probably worth considering.
17
106
Series
Asymptotic Approximation
10
Error
10-4
10-9
10-14
0 20 40 60 100
80
Number of Terms
Figure 5: The relative difference, in absolute value, between the exact value of the function f (ε) = J0 (1/ε), for ε = 1/15, and
the series representation in (38). Also shown is the relative difference between the exact value and the asymptotic expansion
given in (39). The values are given as a function of the number of terms used. Remember that this picture is drawn for one
particular value of ε.
The ideas underlying an asymptotic approximation are so natural that they are found in some of the earliest work in
mathematics. The basis of the definitions given here can be traced back to at least Legendre (1825). In his Traite des
Fonctions Elliptiques he considered approximating a function with a series where the error committed using the first n term
is of the same order as the (n + 1)st term. He referred to this as a semiconvergent series. This is a poor choice of words since
an asymptotic series need not converge. Somewhat earlier than this, Laplace (1812) had found use for asymptotic series
to be able to evaluate certain special functions because the standard series expansions converged so slowly. He considered
the fact that they did not converge to be “inconvenient”, an interesting attitude that generally leads to disastrous results.
The most significant contribution to the subject was made by Poincare (1886). He was able to make sense of the numerous
divergent series solutions of differential equations that had been found up to that time by introducing the concept of an
asymptotic expansion. His ideas are the basis of our development. It is interesting that his paper actually marked the
modern beginnings of two areas in mathematics. One is the asymptotic solution of a differential equation and the other is
the study of divergent series. The latter can also be called the theory of summability of a divergent series, and the books by
Ford (1916) and Hardy (1960) are good references for this area. However, it is not surprising if the reader has not heard of
this particular subject, which has since died of natural causes. For those who may be interested in the early history of the
development of asymptotic approximations, the reviews by McHugh (1971) and Schlissel (1977a) should be consulted.
Exponential asymptotics We mention briefly the topical area of exponential asymptotics, though a study of this is too
advanced for an introductory course on asymptotics like this one. As noted above, useful asymptotic series often diverge;
this is associated with a singularity in the function being approximated and it is now understood that the divergence of the
series is a source of information about the function. Often when using asymptotic series, one is happy with a one or two
term expansion, truncating beyond this point. Occasionally, one requires a higher degree of accuracy. Referring to fig. 5,
which shows the typical behaviour of a divergent (asymptotic) series, one notes that the error first decreases as the number of
terms increases but that at a certain point the error starts to increase. Superasymptotics deals with the idea of finding the
optimal truncation point i.e., finding the number of terms N (this usually depends on the value of ε i.e., N = N (ε)) at which
adding more terms no longer improves the approximation (about N = 47 terms in the figure). It is then sometimes possible
to sum the error if one truncates beyond this point and to show that it is exponentially small with respect to the leading order
terms in the series. Exponential asymptotics (also called asymptotics beyond all orders or hyperasymptotics) goes even
further than this and involves analysing the divergent tail of the series, resumming this, and finding further exponentially
small terms in the function being expanded.
2.6 Stokes phenomenon in the complex plane

An analytic function f (ε) has a power series in ε which converges. If an asymptotic function diverges (this usually involves
exponentially small terms), then this must involve some non-analyticity, e.g. an essential singularity, and so must be restricted
to a sector of the complex plane. There is therefore the interesting possibility of a single function possessing several asymptotic
expansions, each restricted to a different sector of the complex plane. This is called a Stokes phenomenon.
18
Note that while erf as defined below, is analytic at every finite point, there is a non-analytic singularity at ∞.
Z
2 z −t2
erf(z) ≡ e dt.
π 0
2
It can be written as a power series by using the fact that e−t is analytic in the entire complex plane, i.e., it can be expanded
in a Taylor series which converges to the correct value with an infinite radius of convergence. Hence one can integrate term
by term to form a series for erf which also converges with an infinite radius of convergence:

2 1 1
erf z = √ z − z3 + z5 . . .
π 3 10
For z → ∞ with z real, we have the asymptotic expression (obtained by repeated integration by parts,(see Hinch p.20):
2
e−z
erf z ∼ 1 − √ .
z π
2
If z is not real, then we consider the complex plane, and the contour for the integral e−t can be deformed to any complex z
2
with no change in the result so long as the contour is kept in the sector e−z as z → ∞. Thus the above result is applicable
to the sector | arg z| < π4 . A similar expression for the quarter plane about the negative real axis follows from the fact that
erf(z) is an odd function of z:
2
e−z
erf z ∼ −1 − √
z π
as z → ∞ with 3π 5π
4 < arg z < 4 . An expression for erf which is valid at the top and bottom quarter planes, where erf is
very large, can be found by evaluating the original integral from 0 to z using Laplace’s method (see later). The result is:
2
e−z
erf z ∼ − √
z π
as z → ∞ with π4 < arg z < 3π 5π 7π

4 or 4 < arg z < 4 . Thus we have different approximations to the error function which are
applicable to different sectors of the complex plane. Finally we note that the different sectors overlap and do not need to be
mutually exclusive as in the example above.
2.6.1 Manipulating Asymptotic Expansions

It is not hard to show that two asymptotic expansions can be added together term by term. Multiplication is also relatively
straightforward although more tedious and limited to asymptotic sequences that can be ordered in particular ways (e.g., see
Exercises). What is not completely clear is is whether or not an asymptotic approximation can be differentiated (assuming
the functions involved are smooth). To be more specific, suppose we know that
f (x, ε) ∼ a1 (x)φ1 (x, ε) + a2 (x)φ2 (x, ε) as ε ↓ ε0 . (42)

This is sometimes referred to as a generalised asymptotic series because the φi are not just dependent on ε. What is of
interest is whether or not it is true that
d d d
f (x, ε) ∼ a1 (x)φ1 (x, ε) + a2 (x)φ2 (x, ε) as ε ↓ ε0 . (43)
dx dx dx
Based on the conjecture that things will go wrong if at all possible, the only possible answer to the above question is that,
no, it is not always true. To give an example, let f (x, ε) = e−x/ε sin(ex/ε ). For small ε one finds that, for 0 < x < 1,
f ∼ 0 + 0·ε + 0·ε2 . . . . (44)

However, the function
d 1 1
f (x, ε) = − e−x/ε sin(ex/ε ) + cos(ex/ε)
dx ε ε
does not even have an expansion in terms of the scale functions used in (44).
Another problem is that even if φ1 and φ2 are well ordered their derivatives may not be. For example, letting φ1 = 1 + x
d d
and φ2 = ε sin(x/ε), for 0 ≤ x ≤ 1, then φ2 = o(φ1 ) for ε small, but dx φ1 and dx φ2 are not well ordered. So, the question
that we need to address is when can we differentiate an expansion? Well, the answer to this is that if
f (x, ε) ∼ a1 (x)φ1 (ε) + a2 (x)φ2 (ε) as ε ↓ ε0 , (45)

and if
19
1.4
1.2
1.0
5
0.8 y=x -1
0.6
0.4
y
5
0.2
y=x -0.1x-1
-0.0
0.8 0.9 1.0 1.1 1.2
-0.2
-0.4
-0.6 x
Figure 6: Graphs of x5 − 1 and x5 − 0.1x − 1. Both graphs clearly have a root near x = 1.
d
f (x, ε) ∼ b1 (x)φ1 (ε) + b2 (x)φ2 (ε) as ε ↓ ε0 , (46)
dx
d d
then bk = dx ak (i.e., the expansion for dx f can be obtained by differentiating the expansion of f (x, ε) term by term).
Throughout these notes, given (45) we will automatically assume that (46) also holds. Note that a series of the form (45) is
said to be in Poincare form: the φi are dependent only on ε.
The other operation we will have use for is integration, and in this case the situation is better. If the expansion of the
function is as given in (45), and all the functions involved are integrable, then
Z b Z b Z b
f (x, ε)dx ∼ a1 (x)dxφ1 (ε) + a2 (x)dxφ2 (ε) as ε ↓ ε0 . (47)
a a a
It goes without saying here that the interval a ≤ x ≤ b must be within the set of xs where (45) holds.
3 Algebraic and Transcendental Equations

One of the easiest ways to illustrate how an asymptotic expansion can be used to find an approximate solution is to consider
(non-linear) algebraic and transcendental equations. If we consider an equation of the form f (x) = 0 where f is a non-
linear function of its argument, such an equation is either algebraic or transcendental. (An algebraic equation of degree
n is one which can be written as a polynomial of degree n in x where the coefficients in the polynomial are integers e.g.
an xn + an−1 xn−1 + . . . a0 x0 = 0 where each ai ∈ Z. An equation which cannot be written in this form is said to be
transcendental). For our purposes, we will treat algebraic and transcendental equations simply as non-linear equations.
Occasionally the extra information which is available with respect to algebraic equations will be useful e.g. an algebraic
equation of degree n has n roots.
As an aside, note that number theorists classify numbers as being√algebraic or transcendental depending on whether they
can be written as the solution of an algebraic equation or not. Thus, 2 is algebraic (x2 − 2 = 0), but π, e are transcendental.
3.1 Example
To fix ideas first consider the simple example
x5 − 0.1x − 1 = 0. (48)
How might we go about solving this equation? As the equation is a non-linear algebraic equation, in general it will have no
exact solution (unless it’s quadratic or cubic). So we must find an approximation of some sort. The first thing we might do
is graph x5 − 0.1x − 1 against x in which case we get fig.6 and roots occur wherever the graph hits the axis. We see that
there is a root near x = 1 and using a numerical technique (e.g., bisection or Newton-Raphson) we find that x ≈ 1.019607.
In fact we can establish that all the other roots are complex. Let us instead use a perturbation approach. First rewrite the
equation as
20
x5 − εx − 1 = 0 (ε = 0.1) (49)
and a first guess might involve assuming that we can neglect the term εx so that x5 − 1 ≈ 0 and obviously x ≈ 1 is an
approximate solution. Can we do any better? How about trying
x ∼ 1 + εx1 (50)
2
with x1 to be determined. Let us substitute this into (49) and neglect all terms of O(ε ). Thus we find that
(1 + εx1 )5 − ε(1 + εx1 ) − 1 = 0 (51)

and so
1 + 5εx1 − ε − 1 + O(ε2 ) = 0 (52)

and comparing terms of O(ε) we see that x1 = 1/5 and we have the approximation:
x ∼ 1 + ε/5 (53)
For the case ε = 0.1 this gives the approximation x ≈ 1.02 which compares very well with the “exact” numerical solution
x ≈ 1.019607. Furthermore, the approximation holds for all values of ε (provided ε ≪ 1) whereas the numerical solution has
to be recalculated for every new value of ε. In addition, the perturbation solution gives us some feeling for the behaviour of
the solution. It can be improved by going to higher orders and it improves for smaller values of ε. For example if ε = 0.01,
the perturbation solution yields x ≈ 1.002 while the numerical solution is x ≈ 1.001996.
3.2 Example
As another example, consider the quadratic equation
x2 + 0.002x − 1 = 0. (54)
The fact that the coefficient of the linear term in x is much smaller than the other coefficients can be used to find an
approximate solution. To do this, we start with the related question
x2 + 2εx − 1 = 0, (55)
where ε ≪ 1, The beginning assumption is that the solution has a particular form of expansion, namely
x ∼ x0 + εα x1 + . . . , (56)
where α > 0 (this inequality is imposed so the expansion is well ordered for small ε). It should be pointed out that (56)
is nothing more than an educated guess. The motivation for making this assumption comes from the observation that in
expanding functions one usually ends up using Taylor’s theorem and (56) is simply a reflection of that type of expansion.
The exponent α is to allow for a little flexibility.
Now, substituting (56) into (55), one finds that
x20 + 2εα x0 x1 + . . . + 2ε(x0 + εα x1 + . . .) − 1 = 0. (57)

1 2 3 4
The first problem to solve comes from the terms labelled 1 , 4 This is because this equation is supposed to hold for
small ε and therefore we require it to hold as ε ↓ 0. This gives us the following:
ε0 x20 − 1 = 0.
The solutions are x0 = ±1. An important observation here is that this problem has as many solutions as the original
equation (i.e., both problems are quadratic).
With the solution just given, we now have a term in (57) that is O(ε), namely term 3 . Since the right-hand side is zero,
it must be that 3 balances with one of the other terms on the left-hand side of the equation. The only choice is 2 , and
this means α = 1 (Exercise). With this we have the next order problem to solve:
ε1 2x0 x1 + 2x0 = 0
The solution is x1 = −1.
The process used to find x0 and x1 can be continued to systematically construct the other terms of the asymptotic
expansion. We will not pursue this here but point out that the approximations we have calculated so far are x ∼ ±1 − ε. To
illustrate the accuracy of this result, one of the exact solutions is plotted in Fig. 7 along with its asymptotic approximation.
It is clear that they are in very good agreement for small values of ε. Also, if you recall, the original problem was to
construct approximations of the solutions of (54). In this case ε = 10−3 , and so, from the asymptotic approximation we get
x ≈ −1.001, 0.999 while the exact solution gives x = −1.0010005 · · · , 0.9990005 · · · .
21
1.0
0.5
Exact solution
Asymptotic Expansion
0
10-3 10-2 10-1 1
ε - axis
Figure 7: Comparison between the positive root of (55) and the two term expansion x ∼ 1 + ε. It is seen that for small ε the
asymptotic approximation is very close to the exact value. (Note that the horizontal axis is logarithmic).
Note It is appropriate to comment on the notation used in the previous paragraph. We referred to the equation that x0
satisfies as the ε0 equation. It is common to also refer to this as the O(1) equation. However, strictly speaking, based on
the definition of the order symbol, every term in (57) is O(1). In particular we recall that strictly speaking f = o(1) implies
f = O(1). So what we should have done was first introduce a new notation, such as OS (1) (see Exercises). When stating
that f = OS (1), we mean that f = O(1) but f 6= o(1). We could then specify, without ambiguity, exactly what terms go into
the x0 equation. A few textbooks do exactly this to try to reduce the confusion for those who are first learning the subject.
However, the truth is that few people actually use this symbolism. This means O(·) has two meanings, one connected with
boundedness as expressed in the original definition and the other as an identifier of particular terms in an equation.
3.3 Example (A singular perturbation)

Consider the quadratic equation
εx2 + 2x − 1 = 0. (58)
Note that if ε = 0 then the equation becomes linear; that is, the order of the equation is reduced. This is significant because
it fundamentally alters the nature of the equation and can give rise to what is known as a singular problem. Nevertheless,
we will approach this problem in the same way as in the previous example. In other words, we will try the regular expansion
given in (56). Carrying out the calculations, one finds that
1 ε
xa ∼ − + ··· . (59)
2 8
The criticism with this is that there are two solutions of (58) and the expansion has only produced one. One remedy is to
use (59) to factor the quadratic equation (58) to find the second solution. However, there is a more direct way to find the
solution, and it can be adapted to solving differential equations with a similar complication. To explain what this is, note
that the problem is singular in the sense that if ε = 0 then the equation is linear rather than quadratic.
3.3.1 Rescaling the equation

The assumption implicit in (59) (which we assume to be in well scaled non-dimensional form) is that εx2 ≪ 1 and we can
see that this is true a posteriori by examining xa . As we have “lost” a root (x = xb ), there must be something wrong with
this assumption (when we are dealing with x = xb ). Thus for the missing root, εx2 must not in fact be small and we need
some mathematical procedure for dealing with this. In dealing with (58), we assumed that the equation was well scaled i.e.,
that x = O(1) so the terms are of order O(ε), O(1), O(1) respectively; these are the particular balances in the equation which
are appropriate for the first root xa . When searching for the other root, we need to rescale the equation to explicitly reflect
the correct balances associated with the second root. When dealing with algebraic or transcendental equations as here, the
only variable available for rescaling is the unknown quantity x. We will see later when dealing with differential equations
that there are both dependent and independent variables, and on occasion both may need rescaling. Thus we set:
x = εν X (60)
22
where X is a new O(1) variable and so x = Os (εν ) with ν to be determined via re-examining the equation. The equation in
terms of the new variable is:
ε1+2ν X 2 + 2εν X − 1 = 0. (61)

1 2 3
The terms on the left-hand side must balance to produce zero, and we need to determine the order of the problem that
comes from this balancing. There are three possibilities. One of these occurs when ν = 0, and in this case the balance
is between terms 2 and 3 . However, this leads to the expansion in (59) and we have introduced (60) to find the other
solution of the problem. So, we have the following other possibilities:
i) 1 ∼ 3 and 2 is higher order =⇒ 1 + 2ν = 0 =⇒ ν = − 21 =⇒ 1 , 3 = O(1) but 2 = O(ε−1/2 ). This violates our
assumption that 2 is higher order (i.e., 2 ≪ 1 ), so this case is not possible.
ii) 1 ∼ 2 and 3 is higher order=⇒ 1 + 2ν = ν =⇒ ν = −1 =⇒ 1 , 2 = O(ε−1 ) and 3 = O(1). In this case, the
conclusion is consistent with the original assumptions, and so, this is the balancing we are looking for.
With ν = −1, using (61), the rescaled problem is:
X 2 + 2X − ε = 0. (62)
and we note how this displays a different asymptotic balance in comparison with the original equation. We now look for a
straightforward solution of this equation in the form:
X ∼ X0 + εα X1 + . . . , (63)
where α > 0 (so the expansion is well ordered). Substituting this into (62), we have that
(X02 + 2εα X0 X1 + . . .) + 2(X0 + εα X1 + . . .) − ε = 0. (64)

From this, we have:
ε0 X02 + 2X0 = 0 =⇒ X0 = 0, −2 and α = 1

1
ε1 2X0 X1 + 2X1 − 1 = 0 =⇒ X1 = .
2(1 + X0 )
Thus, a two-term expansion of the second solution of (58) is
ε
X ∼ (−2 − ).
2
or in terms of the original variable x we have:
X 1 ε
xb ≡ ∼ (−2 − ). (65)
ε ε 2
Note that the case of X0 = 0 found in the ε0 problem corresponds to the solution given in (59).
3.4 Example
One of the comforts in solving algebraic equations is that we have a very good idea of how many solutions to expect. (Recall
the fundamental theorem of algebra: counting multiplicities, an equation of degree n has n roots (some of which may be
complex of course)). With transcendental equations, this is harder to determine. Perhaps the simplest way to get a handle
on the number of solutions is to graph the functions for specific values of the parameters. To illustrate this idea, consider
the equation
x2 + eεx = 5. (66)
Setting f (x) = x2 + eεx we have that f > 0 and f ′′ > 0. Also f (±∞) = ∞ and f (0) = 1. From this information, we can
make a sketch of f (x), which shows that there are exactly two solutions of this equation. To find asymptotic expansions of
them, assume
x ∼ x0 + εα x1 + . . . .
Substituting this into (66), and using Taylor’s theorem on the exponential (eεx ∼ 1+εx +O(ε2 )), we obtain
x20 + 2x0 x1 εα + . . . + 1 + εx0 + . . . = 5.

From the O(1) equation, we get that x0 = ±2 and α = 1, while the O(ε) equation yields x1 = −1/2. Hence, a two-term
expansion of each solution is x ∼ ±2 − ε/2.
23
10
–3 –2 –1 0 1 2 3
Figure 8: Graphs of y = x2 + exp(εx) and y = 5 for ε = 0.1.
3.5 Example
Find a two term approximation for the solution x of the equation εx tan x = 1 which is near π/2.
To orientate oneself here, we draw a graph of tan x, recalling that the graph is periodic with period π. Note also that
tan π/2 = ∞. We seek a root near π/2 and hence the root must be O(1) so we guess that x ∼ x0 + εx1 .
1
Rearranging the equation we note that tan x = εx and for x = O(1), tan x must be large. For ε → 0, the leading order
approximation is just x ∼ π/2.
Thus we try x ∼ x0 + εx1 = π/2 + εx1 and substitute this into the equation to get:
ε(π/2 + εx1 ) tan (π/2 + εx1 ) ∼ 1 + O(ε2 ).

To progress, we must expand the tan term in powers of ε. We note that:
cos εx1 1
tan(π/2 + εx1 ) = − ∼− + O(1).
sin εx1 εx1
Hence the equation is:
1 π/2
ε(π/2 + εx1 )(− ) ∼ 1 =⇒ − + O(ε) = 1 =⇒ x1 = −π/2
εx1 x1
so our two term solution is x ∼ π/2 − ε π2 .
3.6 Example
We investigate the equation
x
x + 1 + ε sech = 0. (67)
ε
Recall first the identities cosh(x) ≡ 21 (exp(x) + exp(−x)), sinh(x) ≡ 1
2 (exp(x) − exp(−x)) (see figure 9). Note that cosh x
and sinh x are the even and odd parts of exp x respectively.
We assume that ε ≪ 1 is known and we wish to find an approximate solution x for this equation. Sketching f = x + 1 and
g = −ε sech(x/ε) shows that there is one solution of this equation. To find an approximation of it, suppose we proceed in
the usual manner and assume
x ∼ x0 + εα x1 + . . . . (68)
24
7
-0 exp x
-1 cosh x
-2
sinh x
-3
-2 -1 -0 1 2
Figure 9: The exponential functions.
Substituting this into (67), and remembering 0 <sech x≤ 1, it follows that x0 = −1. The complication is that it is not
possible to find a value of α so that the other terms in (67) balance. Putting x = −1 + εα x1 we substitute into (67) to get

α −1 + εα x1
−1 + ε x1 + 1 + εsech
ε
and we need this expression to be small if x is to be an approximate root (if x were the exact root, this expression would
of course be zero). Note that it appears that the leading order terms in the equation balance exactly 1, −1. How about the
next order? Recall that the binomial expansion tells us that:
1 1 1 3 1
(e ε + e− ε )−1 ∼ e− ε − e− ε ∼ e− ε (69)
so
1 1 2 2 1
sech(− ) ≡ ≡ ∼ = 2 exp(− ), (70)
ε cosh (− 1ε ) exp( 1ε ) + exp(− ε1 ) exp( 1ε ) ε
it is clear that:
−1 + εα −1 1
εsech ∼ εsech ∼ 2ε exp −
ε ε ε
1
and this is exponentially small i.e., e− ε ≪ εn for all values of n. This means that the current approach does not give an
approximate root at all as the second and fourth terms in (3.6) cannot balance for any value of α. Most of asymptotics is
concerned with determining the basic balances in equations in the limit of some parameter. (Here, ε → 0). In other words,
our assumption concerning the structure of the second term in (68) is incorrect. Given the character of the hyperbolic
function, it is necessary to modify the expansion and we now assume
x ∼ −1 + µ(ε), (71)
where we are not certain what µ is other than µ ≪ 1 (so the expansion is well ordered). Substituting this into (67), we get
µ + εsech(−ε−1 + µ/ε) = 0. (72)

−1 −1
Now sech(−ε + µ/ε) ∼sech(−ε ) because 1/ε ≫ µ/ε. Recalling again that
1 1 2 2 1
sech( ) ≡ 1/ cosh ≡ ∼ = 2 exp(− ), (73)
ε ε exp( 1ε ) + exp(− 1ε ) exp( 1ε ) ε
because exp( 1ε ) ≫ exp(− 1ε ). (In fact the term exp( 1ε ) is exponentially large). We therefore must have that µ = −2ε exp(− 1ε ).
Thus the correction to the leading order approximation is exponentially small: this explains why the naive approximation
(68) did not work. To construct the third term in the expansion, we would extend (71) and write
25
1
0.8
0.6
0.4
0.2
–3 –2 –1 0 1 2 3
Figure 10: Graph of y =sech x.
Figure 11: The graph of y = x + 1 + εsech ( xε ).
x ∼ −1 − 2εe−1/ε + ν(ε), (74)

where ν ≪ ε exp(−1/ε). To find ν we proceed as before and we need to find the next term in the expansion of sech 1/ε; the
3
details are left as an exercise. (Answer: ν = 2εe− ε ).
One point of this example is to demonstrate that different gauge functions may be appropriate in different situations. Most of
26
the problems that we have dealt with so far have been amenable to the gauge function set {1, ε2 , ε3 ...} but we must be aware
that this is just one particular set and it may not always work. To further illustrate the point, suppose we had approached
the above problem by somehow guessing that the appropriate gauge function set was {1, ε exp(−1/ε)...} and we had thus
written
x ∼ x0 + x1 εe−1/ε + .. (75)
then we would quickly have found that x0 = −1, x1 = −2. However it was not initially obvious that (75) is appropriate; this
is why we introduced (71) and deduced µ(ε) from the analysis.
Again we emphasize that the motivation behind the different approaches is to determine the important asymptotic
balances in the equations.
Large or small roots On occasion one is asked to find the large or small roots or solutions of an equation containing no
small or large parameters. In such problems, the associated small (or large) parameter must be determined as part of the
solution. Thus if one knows that an equation has a large root, x, one might begin by letting x ∼ 1/ε, with ε ≪ 1 to be
determined as part of the solution.
3.7 Iterative methods

An alternative to the expansion methods described above is to from the basis of an iterative (or successive approximations)
process. This can be particularly useful if we do not know the appropriate gauge functions for the problem.
We illustrate the concept with a simple example (which can easily be solved using an expansion process). Consider:
x2 + εx − 1 = 0. (76)
First we note that the exact solutions are: r

ε ε2
x=− ± 1+ . (77)
2 4
Expanding using the binomial theorem gives:
ε ε2 ε4 ε ε2 ε4
x=1− + − + · · · and − 1 − − + + ··· (78)
2 8 128 2 8 128
The first step in the ierative method is to find a suitable re-arrangement of the original equation. Recall that equations of
the form
x = f (x)
can be solved approximately using an iteration scheme (for n ≥ 0):
xn+1 = f (xn ),
for some sufficiently accurate guess x0 . Such a scheme will converge to the root of (76) provided that |f ′ (x)| < 1 for all x
close to the root. First we rearrange the equation (there are many ways to do this: this may require some trial and error):
√
x = ± 1 − εx. (79)
Now if we have an approximation to the root, xn , a better approximation is given by:

√
xn+1 = ± 1 − εxn . (80)
We need a starting point for the iteration: the solution when ε = 0, x0 = 1. After one iteration (on the positive root) we
have: √
x1 = 1 − ε. (81)
If we expand this as a binomial series we find
ε ε2 ε3
x1 = 1 − − − + ··· . (82)
2 8 16
We see this is correct up to ε, but the ε2 terms and higher are wrong. Hence we only retain the first two terms:
ε
x1 = 1 − + ··· . (83)
2
Using this in the next iteration we have: r ε
x2 = 1 − ε 1 − , (84)
2
27
which again can be expanded to give:
ε ε ε2 ε 2 ε3 ε 3
x2 = 1− 1− − 1− − 1− + ··· (85)
2 2 8 2 16 2
ε ε2 ε3
= 1− + + + ··· . (86)
2 8 16
Now the ε2 term is right, but the ε3 terms in still wrong. At each iteration more terms are correct, but more and more work
is required. We can only check to see that a term is correct (without the exact solution) by proceeding to one more iteration
and seeing if it changes.
4 Regular Perturbation Solution of Differential Equations

4.1 Introduction
The ideas used to construct asymptotic approximations of the solutions of algebraic and transcendental equations can also
be applied to differential equations. In this section we consider regular perturbation problems, which roughly speaking,
are problems that retain their essential character as ε ↓ 0. For example, an equation that is second order remains second
order in the limit of ε going to zero. In later chapters, we will take up the study of nonregular, or singular perturbation
problems. We illustrate the main ideas using some examples. A more precise definition of a singular perturbation problem
is that it is an equation which requires more than one scaling.
4.1.1 Example: Projectiles

To illustrate how to find an asymptotic approximation of a solution of a differential equation, we return to the projectile
example described in Section 1.1. Specifically, we consider the problem
d2 y 1
=− , for 0 < τ, ε ≪ 1, (87)
dτ 2 (1 + εy)2
where
y(0) = 0 and y ′ (0) = 1. (88)

The equation in (87) is said to be weakly nonlinear because it is nonlinear but reduces to one that is linear when ε = 0.
To construct an approximation of the solution, we start with a working hypothesis on what is the appropriate form of the
asymptotic expansion. For the problem at hand, we assume
y(τ ) ∼ y0 (τ ) + εα y1 (τ ) + . . . , (89)
where α > 0. As in (56), the exponent α is included to allow for a little flexibility. We will also assume that the expansions
for the derivatives of y(τ ) can be obtained by differentiating (89).
To reduce the differential equation, recall that, for small z, (1 + z)−2 ∼ 1 − 2z. So, substituting (89) into (87) yields
1
y0′′ (τ ) + εα y1′′ (τ ) + . . . = − 2 (90)
[1 + ε(y0 (τ ) + . . .)]
∼ −1 + 2εy0 (τ ) + . . . ,
using the binomial expansion to expand the RHS in powers of ε. From (88),
y0 (0) + εα y1 (0) + . . . = 0 and y0′ (0) + εα y1′ (0) + . . . = 1. (91)

N.B. Note that we must also expand the boundary conditions in powers of ε.
The procedure for finding the terms in the expansion follows the argument used to solve algebraic equations. By equating
like powers of ε, we obtain the following problems:
ε0 y0′′ (τ ) = −1
y0 (0) = 0, y0′ (0) = 1.

The solution of this problem is y0 (τ ) = − 21 τ 2 + τ. Also, because of the second term on the right-hand side of (90), we
must take α = 1 in (89) to get a sensible balance.
ε1 y1′′ (τ ) = 2y0 (τ )
28
0.6
0.4
Solution
Numerical Solution
y0
0.2
y0 + ε y1
0 0.5 1.0 1.5 2.0

τ - axis
Figure 12: Comparison between the asymptotic solutions in (92) and the numerical solution of (87) with ε = 0.1. There is
little difference between the numerical solution and the two term expansion.
y1 (0) = 0, y0′ (0) = 0.

The solution here is y1 (τ ) = 31 τ 3 − 1 4
12 τ .
Therefore, we have that

1 1 1
y(τ ) ∼ τ 1− τ + ετ 3 1 − τ . (92)
2 3 4
This approximation applies for 0 ≤ τ ≤ τh , where τh > 0 is the point where y(τh ) = 0 (i.e., the time when the projectile
returns to the surface of the Earth). A comparison between the asymptotic expansion and the numerical solution of the
problem is given in Fig. 12. It is seen that the two-term expansion is essentially indistinguishable from the numerical
solution, which indicates the accuracy of the approximation.
The second term (y1 ) in (92) contains the contribution of the non-linearity, and the first term (y0 ) is the solution one
obtains in a uniform gravitational field (it is the scaled version of the solution given in (13)). Since y1′ ≥ 0 for 0 ≤ τ ≤ τh
it follows that the contribution of y1 grows with time and it increases the flight time (i.e., it makes τh larger). This is in
agreement with the physics, which says that the gravitational force decreases with height, which should allow the projectile
to stay up longer (Exercise). This is a fairly simple example so not a lot can be said about the solution. What this does
illustrate, though, is that an asymptotic approximation is capable of furnishing insights into the properties of the solution,
which allows us to develop an understanding of the situation that may not be possible otherwise.
4.1.2 Example: A PDE problem

The ideas developed in the last problem can also be applied to partial differential equations. An interesting example arises
in the theory for the diffusion of ions through a solution containing charged molecules (i.e., a polyelectrolyte solution). If the
solution occupies a domain Ω, then the electrostatic potential Φ(x) in the solution satisfies the Poisson-Boltzmann equation
k
X
∇ Φ = − αi zi e−zi Φ ,
2
for x ∈ Ω, (93)
i=1
where the αi s are known positive constants, zi are the known valence of the ionic species, and the sum is over all the ionic
species that are present. For example, if a simple salt such as NaCl is present, then in the sum k = 2 with z1 = −1 and
z2 = 1. There is also an electroneutrality condition for the problem which states that the valences and constants must satisfy
k
X
αi zi = 0. (94)
i=1
On the boundary ∂Ω a uniform charge is imposed, and this condition takes the form
∂n Φ = −ε, on ∂Ω, (95)

where n is the unit outward normal to ∂Ω (note ∂n ≡ n·∇ is the normal derivative).
29
Note For cultural reasons it’s worth pointing out two important special cases of (93). The first arises when the sum contains
two terms, with z1 = −1 and z2 = 1. Using (94) allows (93) to be rewritten as ∇2 Φ = α sinh(Φ), which is known as the sinh
- Poisson equation (or the sinh - Gordon equation). The second special case occurs when there is only one term in the sum
and z1 = 1. In this situation (93) reduces to ∇2 Φ = αeΦ , which is Liouville’s equation. It is also interesting to note that (93)
arises in several diverse areas of applied mathematics. For example, it occurs in the modeling of semiconductors (Markowich
et al., 1990), gel diffusion (Holmes,1990), combustion (Kapila,1983), relativistic field theory (D’Hoker and Jackiw, 1982), and
hydrodynamic stability (Stuart, 1971). It has also been dealt with extensively by MacGillivray (1972) in deriving Manning’s
condensation theory.
Solution The problem we must solve consists of the nonlinear partial differential equation in (93) subject to the boundary
condition in (95). This is a very difficult problem, and a closed-form solution is beyond the range of current mathematical
methods. To deal with this, in the Debye-Hückel theory presented in most introductory physical chemistry textbooks, it is
assumed that the potential is small enough that the Poisson-Boltzmann equation can be linearized. This reduction is usually
done heuristically, but here we want to carry out the calculations more systematically. For our problem, a small potential
means ε is small. Given the boundary condition in (95), it is clear that if ε = 0 the problem becomes:
k
X
∇2 Φ = − αi zi e−zi Φ , for x ∈ Ω, ∂n Φ = 0, on ∂Ω, (96)
i=1
and the solution of the problem is just Φ = 0 as this satisfies the equation and boundary conditions by virtue of the given
condition (94). This suggests that for small ε, the solution will also be small; in fact we expect it to be O(ε). We like to deal
with scaled quantities which are O(1) at leading order. Thus it is sensible at this point to rescale the problem to reflect this
fact. So we write:
Φ = εφ (97)
where φ is the new O(1) variable so that Φ = O(ε) when φ = O(1). Now the problem becomes:
k
X
ε∇2 φ = − αi zi e−zi εφ , for x ∈ Ω, (98)
i=1
with rescaled boundary condition:
∂n φ = −1, on ∂Ω. (99)

The appropriate expansion of the potential, for small ε, is
φ ∼ φ0 (x) + εφ1 (x) + . . . . (100)

Substituting this into (93), we obtain
ε(∇2 φ0 + ε∇2 φ1 + . . .) (101a)

X
= − αi zi e−εzi (φ0 +εφ1 +...)
X
1 2 2 2
∼ − αi zi 1 − εzi (φ0 + εφ1 + . . .) + zi ε φ0 + . . .
2
X X X 1
∼− αi zi + ε αi zi2 φ0 + ε2 αi zi2 (φ1 − zi φ20 ) + . . . .
2
From the electroneutrality condition in (94), and setting
k
X
2
κ = αi zi2 ,
i=1
it follows that the (leading order) ε0 problem is
∇2 φ0 = κ2 φ0, for x ∈ Ω, (102)

Expanding the boundary condition (99) we find that
∂n (φ0 + εφ1 + ...) = −1, on ∂Ω (103)

so the ε0 boundary condition is:
30
∂n φ0 = −1, on ∂Ω. (104)
To solve this problem, the domain Ω needs to be specified; an example follows.
The first-term approximation Φ ∼ εφ0 (x) is the basis of the classical Debye-Hückel theory in electrochemistry. One of the
advantages of our approach is that it is relatively easy to find the correction to this approximation, and this is accomplished
by solving the ε1 problem. From (101a), one finds that
(∇2 − κ2 )φ1 = −λφ20 , for x ∈ Ω, (105)

X
1
where λ = 2 αi zi3 and, from (103) the relevant boundary condition is:
∂n φ1 = 0, on ∂Ω. (106)
It is assumed here that λ 6= 0.
4.1.3 Exercise
It remains to solve the problems we have derived for the first two terms in the expansion. In electrochemistry, an often-used
domain in two-dimensions is the region in the plane outside of the unit circle (i.e., the region 1 < r < ∞). Using polar
co-ordinates, the solution is independent of the angular variable and both (102) and (105) reduce to ordinary differential
equations in the radial coordinate i.e., Φ = Φ(r) and the governing equation becomes an ordinary differential equation with
d2 Φ 1 dΦ
∇2 Φ ≡ + (107)
dr2 r dr
with boundary conditions:
dΦ
= −1 on r = 1 (108)
dr
and we require that Φ be finite as r → ∞. Carrying out the calculations, one finds that the solution of (102,104) is
K0 (kr)
φ0 = − ,
kK1 (k)
where K0 and K1 are modified Bessel functions. With this it is not difficult to solve (105) for φ1 using, for example, the
method of variation of parameters.
Our approach to regular perturbation expansions is constructive, and in the development little has been said about the
theoretical foundation of the subject. For example, we have not made any attempt to prove rigorously that the expansions
are asymptotic. This requires estimates of the error, and a nice introduction to this subject can be found in Murdoch (1991).
4.1.4 Perturbed boundary conditions

Suppose that one is dealing with the following problem:
d2 y dy
+ = 1; y(x = ε) = 1, y(x = 1) = 0,
dx2 dx
how does one deal with the perturbed boundary condition? It is clear that the zero order problem can be obtained in the
usual way by just putting ε = 0. But what happens at the next order? The difficulty arises with the perturbation of the
boundary in the sense that the unperturbed left hand boundary is x = 0. Suppose we expand in the usual way:
y ∼ y0 (x) + εy1 (x).

It should be clear that the expansion of the equation and the right hand boundary condition are trivial. The only difficulty
arises when one deals with the perturbed boundary. Specifically if we expand this condition as
y0 (x = ε) + εy1 (x = ε) = 1
the question arises as to the precise order of each term. y0 (x = ε) looks at first glance to be O(1) but it is not of course,
because it still has an implicit dependence on ε. A little reflection suggests that the way to deal with this boundary condition
is merely to expand in a Taylor series about ε = 0. In this case we have
dy0 (0)
y0 (x = 0) + ε + εy1 (x = 0) + O(ε2 ) = 1
dx
and now it is clear that we have successfully segregated the terms in the boundary condition into ε0 and ε1 contributions.
Thus, as regards the boundary condition, we have the following conditions:
ε0 : y0 (x = 0) = 1; ε1 : y1 (x = 0) = − dydx
0 (0)
.
31
4.2 Uniformity of approximations
It is often the case that the functions expanded in an asymptotic expansion depend on more than just the expansion
parameter. This can interfere with the accuracy, or even existence, of an asymptotic expansion. To illustrate the situation,
consider the following algebraic equation for y:
(y − 1)(y − x) + εy = 0.
Suppose we are interested here in values of x in the interval 0 < x < 1. For small ε the appropriate expansion of the solution
is y ∼ y0 + εy1 + . . . . Substituting this into the equation and carrying out the calculations, one finds the following expansion
for one of the solutions:
ε
y ∼1− + ···
1−x
There is no doubt that this is an asymptotic expansion of the solution; the only question is how accurate it is for values
of x that are close to x = 1. Given the nature of the O(ε) term in the expansion, it is expected that it will be necessary to
have very small ε if such x are used. To say this in a slightly different way, suppose we are going to use only the first-term
approximation y ∼ 1 and we want to determine how small ε should be so the error in this approximation is less than, say,
10−4 for all the x′ s under consideration. Using the second term as a measure of the error, we are asking if it is possible to
find an interval 0 < ε < ε2 so ε/(1 − x) < 10−4 for 0 < x < 1. Clearly this is not possible because ε/(1 − x) is unbounded as
a function of x over this interval. For this reason, the expansion is not uniformly valid for 0 < x < 1. On the other hand, it
is uniformly valid if, say, 0 < x < 3/4.
In fact we can be more precise about the region where the above expression fails. Clearly the right hand side is no longer
ε
well ordered when 1 and 1−x are of the same order of magnitude. Fix ε and let x = 1 + δ(ε) where δ(ε) ≪ 1 is to be
determined. For x = 1 + δ(ε) the two terms have magnitude 1; δε and these are of the same order if δ = O(ε) so the expression
fails within an O(ε) region about x = 1.
The following definition is introduced to make precise what is meant by a uniform approximation (c.f. 2.4.1). It is written
for the specific case of small ε because that is the only limit we will be considering when dealing with uniformity.
4.2.1 Definition of uniform asymptotic approximation

Suppose f (x, ε) and φ(x, ε) are continuous functions for x from a set I and ε from an interval 0 < ε < ε1 . In this case, φ(x, ε)
is a uniformly valid asymptotic approximation to f (x, ε) for any x in I as ε ↓ 0 if f − φ = o(φ) in this limit (or equivalently
if φf → 1 ).
The crucial idea here is that we need to consider ε to be small (as small as we like) but fixed and in such circumstances
we ask whether f − φ = o(φ) for all values of x i.e., we are now considering limits with x varying.
4.2.2 Definition of uniform asymptotic expansion or series

The natural extension to this is to say that relative to the well-ordered set of gauge functions φ1 (ε), φ2 (ε)...
∞
X
f (x, ε) ∼ φi (x, ε) = φ1 (x, ε) + φ2 (x, ε) + .. (109)
i=1
is a uniformly valid asymptotic expansion or series of f (x, ε) if each partial sum is a uniformly valid asymptotic approxi-
mation ∀x ∈ I i.e.
(f − φ1 ) = o(φ1 ); (f − φ1 − φ2 ) = o(φ2 ) . . . (110)
It can be shown that such a series is well-ordered i.e., that each successive term is asymptotically smaller than the preceding
one.
As usual, given a function f (x, ε), one of the main tools that is used to establish uniformity is Taylor’s theorem, which
is illustrated in the following examples.
Examples
π 3π
1. If f = sin(x + ε), then f ∼ φ ≡ sin(x) is a uniformly valid asymptotic expression for 4 < x < 4 . This can be
established using Taylor’s theorem (about x = O(1) and ε ≪ 1),which gives us that
f = sin(x) + ε cos(x) + O(ε2 )

f
for all x ∈ ( π4 , π4 ). We notice that all the x functions in this series must be O(1). Thus, limε↓0 φ = 1 for all x ∈ ( π4 , π4 )
and from this we have that f ∼ sin(x) for π4 ≤ x ≤ 3π 4 .
What may be surprising is that it is not a uniform approximation for 0 < x < π2 . To see this, we fix ε ≪ 1 and consider
sin(x + ε). For small x, we have sin(x + ε) ∼ x + ε, while sin(x) ∼ x. In this case lim x+ε
x does not tend to 1 as x → 0
32
with ε fixed. Alternatively we see that it is not true that (x + ε) − x ≪ x for all values of x. No matter how small (but
fixed) we choose ε, there exist x such that (x + ε) − x = O(x) and also there exist x such that (x + ε) − x ≫ x so the
approximation is not uniformly valid.
Thickness of region of non-uniformity More precisely, fix ε and put x = 0 + δ(ε) where δ(ε) ≪ 1 is to be
determined. Now
f = sin(δ + ε) ∼ δ + ε; φ = sin(δ) ∼ δ (111)
f ∼ φ requires that lim δ+ε

δ = 1 and this fails when δ = O(ε) so the approximation φ = sin x fails in an O(ε)
neighbourhood of x = 0.
2. A solution of the equation
εy ′′ + 2y ′ + 2y = 0, for 0 < x < 1, (112)

is
αx βx
y(x) = e , ε +e ε (113)
√ √
where α = −1 + 1 − 2ε and β = −1 − 1 − 2ε. To obtain an asymptotic approximation (not an expansion) of this
solution for small ε, note α ∼ −ε − ε2 and β ∼ −2 + ε. Thus, for 0 < x < 1,
−ε−ε2 −2+ε
x x
y(x, ε) ∼ e ε +e ε (114)
∼ e−x + e−2x/ε (≡ φ1 (x, ε)) (115)
∼ e−x (≡ φ2 (x, ε)). (116)

φ1 is a uniform approximation for y because if we fix ε, then φy1 → 1 uniformly for all x. However, φ2 is not a uniform
approximation for y or for φ1 for 0 < x < 1. This is because the exponential that was dropped in going from (115) to
(116) is of the same order as (116) when x = 0. Thus φφ21 → 1 + ex(1−2/ε) and no matter how small we choose ε, we
can find x = O(ε) such that ex(1−2/ε) is exp(O(1)) = O(1). Similarly we can show that φy2 9 1 uniformly.
Thickness of region of non-uniformity To be specific compare the expressions: e−x + e−2x/ε ; e−x with ε fixed.
Let x = 0 + δ(ε) where δ(ε) ≪ 1 is to be determined. (115) and (116) are now of order
e−δ + e−2δ/ε ; e−δ . (117)
The first and last of the three terms are always O(1) for 0 < x < 1. They will be approximately equal to 1 when
δ = O(ε) because e−ε ∼ 1. The middle term will become exp(−O(1)) when δ = O(ε) and will then begin to make a
significant contribution. Thus (115) and (116) will start to vary significantly when x is within an O(ε) neighbourhood
of 0.
√
It may help to note that when δ is small but not as small as ε, for example δ = ε ≫ ε, then the terms are of size
√ √
− √2ε
e− ε
;e ; e− ε
,
which are O(1), exponentially small and O(1) respectively.

Similarly
y(x) ∼ e−x + e−2x/ε , (118)

is a uniformly valid approximation of (114) over the interval 0 < x < 1 while (116) is not. To illustrate the differences
between these approximations, they are plotted in Fig. 13 in the case of ε = 10−2 . Clearly, the uniform approximation
gives a very good description of the function (114) over the entire interval. Also, over most of the interval (116) is
close to the exact solution (114). The exception, as expected, is the region near x = 0, particularly in the region where
x = O(ε). This region is an example of what is known as a boundary layer and will be investigated extensively in the
next chapter.
33
2.0
Exact Solution
Nonuniform Approximation
1.6
Uniform Approximation
1.2
Solution
0.8
0.4
0
10-3 10-2 10-1 1
x - axis
Figure 13: Comparison between the exact solution (113) of the differential equation, a nonuniform approximation (116), and
a uniform approximation (118). (The uniform approximation and the exact solution are almost indistinguishable). ε = 0.01.
Note the logarithmic scale on the x axis.
Figure 14: Comparison between the exact solution (113) of the differential equation, a nonuniform approximation (116), and
a uniform approximation (118). (The uniform approximation and the exact solution are almost indistinguishable. ε = 0.01.
Here we use a standard linear scale on the x axis.
34
4.3 Small and large time solutions; artificial parameters
On occasion one wishes to have a small or large time solution (or the “near or far field” behaviour i.e., solutions for small
or large x) to a problem. This may be of interest in itself or may be used to start up or to verify a numerical solution of a
problem. Even if there are no other small or large parameters, it is possible to use the smallness or largeness of t to obtain
an asymptotic approximation.
As t or x is a variable quantity, it is simplest to proceed by introducing an artificial parameter. For example, if we are
seeking a small time solution, i.e., a solution for small t, then we set
t = ετ ; ε ≪ 1, τ = O(1). (119)
Obviously for a large time solution we would typically write t = τ /ε, although other scales may be necessary for more
complicated problems. The problem thus becomes a parameter perturbation problem and one proceeds as previously with a
small parameter by scaling all variables to obtain the balances which are appropriate for small t. There is one extra point:
the final solution, when written back in terms of the original variable(s) (i.e., in terms of t not τ ) should be independent of
the artificial small parameter. This is obvious in the sense that if we were to look for a small time expression for sin t, then
3
this is clearly going to depend on t alone. Thus we know that sin t ∼ t − t3! + . . .. But if we had formally written
(ετ )3
t = ετ ; sin (ετ ) = ετ − + ... (120)
3!
then we see that each term in the asymptotic series, if rewritten in terms of t, is independent of ε.
Example Consider the problem (similarity solution of the diffusion equation)
d2 f df
ε2 2
+x = 0; f (0) = 0, f (∞) = 1. (121)
dx dx
If we change variables as follows:
x̄ = x/ε, g(x̄) = f (x) (122)
the problem becomes:
d2 g df
2
+ x̄ = 0; g(0) = 0, g(∞) = 1, (123)
dx̄ dx̄
and the parameter ε disappears completely from the problem. It is thus an artificial parameter.
Obviously the solution of the problem for g(x̄) is of the form g = g(x̄) while the solution of the x problem is of the form
f = f (x, ε). But because we were able to remove ε from the second formulation it must be possible to write the solution of
the first (x) formulation independently of ε and each term of this solution must be capable of being written independently of
ε when it is rewritten in terms of x̄. Of course, for consistency with the rest of this section, we should probably think of the
x̄ formulation as comprising the original problem and the x formulation as the version incorporating the artificial parameter
ε. To verify what we have just said, as an exercise check that the solution of the problem is:
x x
f = erf ( √ ) =⇒ g = erf ( √ ).
2ε 2
In some cases, small time solutions can be found by using Taylor series. The small parameter method comes into its
own when the solution is not real analytic (i.e., does not have a Taylor series so it does not have a representation in positive
integer powers of t).
Example Find the small t behaviour of the ODE

dy 1
+ ty = 4 . (124)
dt t
We write:
dx 1
t = ετ, y = εa x =⇒ εa−1 + εa+1 τ x = ε−4 4 , (125)
dτ τ
where ε ≪ 1 is an artificial small parameter and a is to be determined by balancing terms. We have three terms which are
O(εa−1 ), O(εa+1 ) and O(ε−4 ) respectively. If we balance the first two terms and assume the third is much smaller then we
need to choose a − 1=a + 1 which is impossible. If we balance the second and third terms (assuming the first is much smaller)
we require a + 1=−4 and so a = −5 but then the first term is O(ε−6 ) and it dominates for ε ≪ 1. Finally if we balance the
first and third terms then a − 1=−4 and so a = −3 and the second term is O(ε−2 ) which is much smaller. Thus a = −3 and
dx 1 1 1
O(ε−4 ) : = 4 =⇒ x = − 3 + B. (126)
dτ τ 3τ
35
Thus we have
1 1 B 11 B
ε−3 x ∼ −ε−3 3
+ 3 =⇒ y ∼ − 3 + 3
3τ ε 3t ε
reverting to the original variables. As we noted above, when we write the solution in terms of the original variables ε should
not appear. Thus we must choose B = 0 and the leading order behaviour is independent of any initial condition (no initial
condition was given here but even if it had been the leading order behaviour does not depend on it). The final result for the
leading order small time behaviour is thus:
11
y ∼ − 3 , t → 0. (127)
3t
We note that the above approach only covers the situation where the small time expansion of the dependent variable y is
in powers of t. Let us dogmatically check this small t solution by substituting (127) into (124) and recalling that we only
require the equation to be satisfied in the limit as t → 0 i.e., we really require that:
dy 1
+ ty ∼ 4 , t → 0.
dt t
On substituting in we find that:
1
t−4 − t−2 ∼ t−4 , t → 0
3
which is certainly true as t−4 ≫ t−2 as t → 0.
It is possible to proceed to higher terms in the usual way by substituting
x ∼ x0 + εb x1 + . . .
into (125) and determining the powers by asymptotic balance and one finds that
11 1 t t2
y∼− 3
− + C + − C . . . , t → 0.
3t 3t 3 2
The constant C could be determined from an initial condition (e.g., y(1) = 1).
Example Find the decaying large time solution of

du
= −u + u2 /2. (128)
dt
For this equation it turns out that there are two particular large time solution behaviours depending on the initial condition.
Because we can solve this problem exactly, we will first cheat and find the exact solution. First we graph du/dt against u and
the corresponding “phase line” (this is just the horizontal axis of figure 15. From the figure we see that there are two fixed
points at u = 0 and u = 2. For initial conditions with −∞ < u(0) < 2 solutions will tend to the zero (decaying solutions)
while if u(0) > 2 the solution will tend to ∞. We can verify this by solving the problem directly. The ODE is first order,
nonlinear, but fortunately the variables are separable. We rewrite it as:
du
= dt (129)
u(u/2 − 1)
and using partial fractions we find that:

1 1 1
=− +
u(u/2 − 1) u u−2
so the solution is:
2
u= ,
1 + 2Aet
2
where A is determined from the initial condition. For example if u(0) = 1 then u(t) = 1+e t and for large t we see that
−t
u ∼ 2e .
Now suppose that we did not have this exact solution. To find the large time decaying solution, we might proceed as
follows. Let ε ≪ 1 be an artificial parameter. We rescale via:
t = f (ε) + τ, u(t) = εv(τ ),
where f (ε) ≫ 1, the time scale for u to decay to O(ε) is to be determined in the course of the analysis. The scaling is
motivated by the fact that for large times we expect that the first two terms in (128) must balance. The rescaled equation
is:
dv 1
= −v + εv 2 . (130)
dτ 2
36
Figure 15: The phase line.
Seeking a series solution v ∼ v0 + εv1 + . . ., we find at leading order that:

dv0
ε0 : = −v0 =⇒ v0 ∼ Be−τ . (131)
dτ
At this point f (ε) has still not been determined. However if we revert to the original variables we find that:
1
u ∼ εBe−t ef (ε ) =⇒ f (ε) = − ln ε ≡ ln ( ), (132)
ε
in order to ensure the solution is independent of the artificial parameter. Note that this means that u will decay to O(ε) in
a time of O(ln 1ε ). Thus our final leading order estimate is:
u ∼ Be−t (t → ∞). (133)
B can be estimated by asymptotic matching with an intermediate time solution or by numerical patching.
An alternative approach would have been to scale initially:
t = ε + τ ; u = g(ε)v(τ )
with g(ε) to be determined. (Exercise: show that g(ε) = e−1/ε ).

Sometimes far field or large time solutions can be obtained using a WKB kind of approach (see chapter 4) though this
technique is only applicable for linear equations.
Exercise Consider the PDE problem (diffusion or heat flow on x ∈ [0, π]):
ut = uxx , u(x, 0) = sin x, u(x = 0, t) = u(x = π, t) = 0. (134)
Verify that this has the exact solution u = sin xe−t and expand this for small t to get u ∼ sin x(1 + t + t2 /2 + . . .). Now use
the artificial small parameter method to obtain a two term small time solution and compare the results.
37
4.4 Maple
Anyone who uses perturbation methods is struck almost immediately by the amount of algebra that is sometimes necessary
to find an approximation. Because one of the objectives of the derivation is to have an answer that contains the parameters
in the problem, it is not convenient to use a computing language like FORTRAN or C to carry out these calculations since
they require all variables (and parameters) to be evaluated. An alternative is to use a symbolic language like Macsyma,
Maple, or Mathematica. Symbolic manipulators have been around for some time but only recently have they become widely
used. For those who may be unfamiliar with these systems, they are capable of manipulating mathematical expressions
without having to evaluate parameters. This ability has a tremendous potential for constructing asymptotic expansions, and
so we will illustrate some simple applications. More extensive presentations can be found in Rand and Armbruster (1987)
and Heck (1993).
The simplest example is in the use of Taylor’s theorem. For example, suppose we want a two-term expansion of the
function

1
f (ε) = sinh ,
10α + cos(ε)
where α is a constant. If we expect that the expansion of this function is a regular power series, that is,
f ∼ f0 + εf1 + ε2 f2 + . . . ,
then we can obtain the coefficients using a Taylor expansion. To do this, one must enter the formula for f, and in the
following command ep designates ε and a represents α. In Maple, the command for the Taylor series expansion is then
taylor(f, ep, 3); (135)

In a fraction of a second the machine dutifully responds with
 
exp 1
exp − 1
1 1 1 1 1 10a+1 1 10a+1
 ep2 + O(ep3 ).
exp − exp − + +
2 10a + 1 2 10a + 1 4 (10a + 1)2 4 (10a + 1)2
Clearly, the effort involved here is minimal, and this makes the use of a symbolic manipulator quite attractive. However,
at the same time certain cautions are necessary. It is natural to think that since this is so easy then why not go for lots
of terms and change the 3 in (135) to, say, 1000. There may actually be situations where something like this is necessary,
but one would risk creating what is known as intermediate expression swell. This is the situation where the expressions
being manipulated are so large that the machine is incapable of handling them with the allotted memory even though the
final result is relatively simple. In terms of asymptotic approximations, such Herculean feats as deriving 1000 terms of an
expansion are generally of little value. In fact, as we have seen earlier, the results may actually get worse. It should also be
pointed out that even though these systems can be big time savers there are certain things they do not do easily, such as
simplifying certain expressions. For example, the above expansion can be written as

cosh 1
1 10a+1
sinh + 2 ep2 + O ep3 .
10a + 1 2 (10a + 1)
It can be very difficult to get the program to “notice” this fact. Also, the program does not warn the user that this expansion
is singular when 10a + 1 = 0. The point here is that one must not stop thinking about the problem, or its solution, just
because the work is being done by a machine.
The expansions for the solutions of many of the algebraic and differential equations studied in this chapter can be solved
quite easily using a symbolic manipulator. To illustrate, consider the projectile problem (15). The Maple commands that
can be used to construct the expansion are given in the Table below. The cpu time needed to carry out these commands
depends on the machine being used; however, it should not be more than about five or ten seconds. Also, even though the
example presented here uses Maple, it is possible to carry out similar calculations using most other symbolic programs.
It is strongly recommended that anyone using asymptotic expansions experiment with a symbolic computing system using
a program like the one described here.
Command Explanation
restart;ep:=epsilon; resets all variables
y:=y0(t)+ep*y1(t)+epˆ2*y2(t); defines the terms of the expansion
de:=diff(y,t,t)+1/(1+ep*y)ˆ2; defines the differential equation
deser:=taylor(de,ep,3); expands equation in powers of ep
de0:=coeff(deser,ep,0); O(1) equation
de1:=coeff(deser,ep,1); O(ep) equation
y0(t):=rhs(dsolve({de0,y0(0)=0,D(y0)(0)=1},y0(t))); solve O(1) problem +ICs
y1(t):=rhs(dsolve({de1,y1(0)=0,D(y1)(0)=0},y1(t))); solve O(ep) problem
38
Other Maple information
• On beginning choose file-new-document mode.
• Using document mode, ctrl-shift-k inserts new line before current line.
• Using document, ctrl-shift-j inserts new line after current line.
• ctrl-delete deletes current line (do not highlight).
Here is a sketch of some useful code to solve a differential equation numerically and draw a graph of the answer.
Command Explanation
restart; resets all variables
with(plots): gives access to a plot package, colo
y0:=y(t);yd:=diff(y0,t); defines the terms in the equation
ode:=yd+y0-1; defines the differential equation
t1:=dsolve({ode,y(0)=1},y(t),numeric); solve the equation numerically; M
ploty:=odeplot(t1,[t,y(t)],t=0..3); plot the solution and place result
display(ploty); display the solution
# Some other useful commands: with(Typesetting): Settings(functionassign=false); does not assign functions to statem
dchange useful for change of variables and
4.5 The chain rule for partial derivatives

Suppose that F = F (x, y, z) is a continuously differentiable function of x, y, z. Suppose further that
x = x(r, s, t); y = y(r, s, t); z = z(r, s, t) (136)

are continuously differentiable functions of r, s, t. Then F is also a continuously differentiable function of r, s, t. For clarity it
is conventional to use a different label for F when it is considered a function of r, s, t. Thus we write F (x, y, z) and G(r, s, t)
though F and G represent the same function.
The chain rule for evaluating the partial derivatives of G(r, s, t) is then:
∂G ∂F ∂x ∂F ∂y ∂F ∂z
= + + (137)
∂r ∂x ∂r ∂y ∂r ∂z ∂r
∂G ∂F ∂x ∂F ∂y ∂F ∂z
= + +
∂s ∂x ∂s ∂y ∂s ∂z ∂s
∂G ∂F ∂x ∂F ∂y ∂F ∂z
= + + .
∂t ∂x ∂t ∂y ∂t ∂z ∂t
The following important subcase often arises. Suppose that F = F (x, y, z) and x = x(t), y = y(t), z = z(t) with the same
restrictions as above. Then we may consider F to be a function of t which we label G(t) and:
dG ∂F dx ∂F dy ∂F dz
= + + . (138)
dt ∂x dt ∂y dt ∂z dt
Example
Demonstrate the chain rule for the case:
F (x, y, z) = x + y 2 + z 3 ; x = t; y = t2 ; z = t3 . (139)
Solution Using the chain rule, when F is considered a function of t, which we label G(t) then:
dG ∂F dx ∂F dy ∂F dz
= + + = 1 + 4yt + 9z 2 t2 . (140)
dt ∂x dt ∂y dt ∂z dt
It is sensible to write this wholly in terms of t whence we obtain: dG/dt = 1 + 4t3 + 9t8 .
We can also write down the function G explicitly i.e.,
dG
G(t) = t + t4 + t9 =⇒ = 1 + 4t3 + 9t8 (141)
dt
and the chain rule is verified.|
39
∂G ∂G
Exercise Demonstrate the chain rule (i.e., evaluate ∂r , ∂s directly and using the chain rule) for the case: F = x + sin y
and x = r cos s, y = r sin s.
4.6 Solutions to some elementary ODEs.

We take y = y(x) to be the dependent variable and x to be the independent variable.
1. First order: variables separable

This has the basic form:
dy
g(y) = −f (x) (142)
dx
where f (x) is a given function of x and g(y) is a given function of y. If the equation has this form, the variables can
be separated to get all xs on one side and all ys on the other:
Z Z
g(y)dy = −f (x)dx → g(y)dy = − f (x)dx + constant. (143)
The equation can now be solved by quadrature as in principle both of these integrals can be evaluated.
N.B. g(y) can be a non-linear function of y and this method will still work in principle.
Example
(From separation of variables method for PDEs)

dG
Find the general solution of: dt = kG where k is a known constant.
Solution
dG dG dG
= −kG ⇒ = −kdt ⇒ ∫ = −k ∫ dt ⇒ (144)
dt G G
1nG = −kt + C ⇒ G = exp(−kt + C) = exp(C)exp(−kt) = Dexp(−kt) (145)
where C and D are arbitrary constants. Note the step whereby exp(C) is replaced by D. As C is arbitrary at this
point, so is exp(C) and we just call this D. This step is not essential but it puts the final solution in a simpler form.
D is found by using the boundary conditions for the problem.
2. First order linear equation (wrt y)
This equation has the form:
dy/dx + P (x)y = Q(x) (146)
where P (x) and Q(x) are knownR functions. This equation can be solved by multiplying both sides by an integrating
factor which is of the form exp( P (x)dx). (The equation must be in exactly the same form as (146) i.e., coefficient
of dy/dx must be one. This can always be arranged by dividing across the equation by the coefficient of dy/dx.) The
equation is structured
R in such a way that on multiplying across by the integrating factor, the LHS of the equation now
equals d/dx(y exp ( P (x)dx)) and in this way the equation can be integrated directly.
Example
Find the general solution of 2dy/dx + 2y/x = 4x.
40
Solution First divide the equation by 2 to make the coefficient of dy/dx equal. Thus we have the equation dy/dx +
y/x = 2x. Here P (x) ≡ 1/x and Q(x) ≡ 2x. Thus the integrating factor is exp(∫ 1/xdx) = exp (ln(x)) = x. So we
multiple both sides of the equation by x to get:
x(dy/dx + 1/xy) = 2x2 (147)
The LHS now will always equal d/dx(y × integrating factor) i.e., d/dx(yx) in this example. This is the key point behind
the technique. (Check that this is so) and the equation becomes:
d(xy)
= 2x2 ⇒ d(xy) = 2x2 dx ⇒ ∫ d(xy) = 2 ∫ x2 dx ⇒ xy = 2x3 /3 + C ⇒ y = 2x2 /3c + C/x (148)
dx
where C is an arbitrary constant to be determined from the boundary conditions.
3. Linear constant coefficient equations can be solved using the exponential substitution y = ekx where k is a constant
to be determined.
4. Simple harmonic motion equation
d2 y
= −n2 y (149)
dx2
where n is a given real constant.
This has general solution:
y = A cos nx + B sin nx
where A and B are arbitrary constants to be determined from the boundary conditions.
To obtain this solution, we note that the equation is second order, linear (in y), homogeneous and with constant
coefficients. Thus we recall that the exponential substitution y = exp(kx) for some constant k will work in principle.
Substituting in we find that dy/dx = k exp(kx) and d2 y/dx2 = k 2 exp(kx) so the equation reduces to:
k 2 exp(kx) = −n2 exp(kx) ⇒ k 2 = −n2 ⇒ k = ±in

(because n is a real number).
So we have (complex) solutions of the form exp (inx) and exp(−inx). In order to get real solutions, we recall that
if we have a complex solution of a linear homogeneous ODE, both the real and imaginary parts must separately be
solutions. It is easy to see why this is the case. If yc (x) = yr (x) + iyi (x) is a complex solution of L[y]=0 then using
the linearity property
L[yr (x) + iyi (x)] = O + Oi = L[yr (x)] + iL[yr (x) + iyi (x)]
and so we must have L[yr (x)] = O and L[y1 (x)] = O. Now we use Euler’s result that
exp(iθ) = cosθ + i sin θ.
Thus exp (inx) = cos nx = i sin nx and exp (−inx) = cos nx − i sin nx so we see that cos nx and sin nx are fundamental
(and linearly independent) solutions. As the equation is second order and linear the general solution is just A cos nx +
B sin nx for arbitrary constants A, B (to be determined from the boundary conditions of the problem).
5. The equation:
d2 y
= n2 y (150)
dx2
where n is a given real constant. This can be solved using the same technique as above. It has general solution:
y = A cosh nx + B sinh nx (or equivalently C exp (nx)) + D exp (−nx)) where A, B, C, D are arbitrary real constants to
be determined from the boundary conditions.
Recall that cosh and sinh are the even and odd parts of the exponential function exp. Thus cosh(x) ≡ 1/2(exp (x) +
exp (−x)) and sinh(x) ≡ 1/2(exp (x) − exp (−x)) whence cosh(−x) = cosh(x) and sinh(−x) = sinh x.
41
6. Bessel’s equation of order n (linear with non-constant coefficients).
d2 y dy
x2 2
+x + (x2 − n2 )y = 0
dx dx
with general solution y = AJn (x) + BYn (x) where A and B are arbitrary constants to be determined from the boundary
conditions and Jn and Yn are the Bessel functions of the first and second kind of order n. Jn and Yn are known functions
(just as cos and sin are known functions) and values for these functions can be looked up in tables (just as for cosine
and sine).
7. Legendre’s equation of order n (linear with non-constant coefficients).
d2 y dy
(1 − x2 ) − 2x + n(n + 1)y = 0. (151)
dx2 dx
with general solution y = APn (x) + BQn (x) where A and B are arbitrary constants to be determined from the boundary
conditions and Pn and Qn are the Legendre function of the first kind of order n and Legendre function of the second
kind of order n. Pn and Qn are known functions (just as cosine and sine are known functions) and values for these
functions can be looked up in tables (just as for cos and sine).
8. The simplest second order ODE (d2 y/dx2 = 0).
d2 y d dy dy
=0⇒ ( )=0⇒ = C ⇒ dy = C dx ⇒ y = Cx + D (152)
dx2 dx dx dx
5 Matched Asymptotic Expansions

5.1 Introduction
The ideas underlying an asymptotic approximation appeared in the early 1800’s, when there was considerable interest in
developing formulas to evaluate special functions. An example is the expansion of Bessel’s function, given in (39), that was
derived by Poisson in 1823. It was not until later in the century that the concept of an asymptotic solution of a differential
equation took form, and the most significant efforts in this direction were connected with celestial mechanics. The subject of
this chapter, what is traditionally known as matched asymptotic expansions, appeared somewhat later. Its early history is
strongly associated with fluid mechanics, and specifically, aerodynamics. The initial development of the subject is credited
to Prandtl (1905), who was concerned with the flow of a fluid past a solid body (such as an airplane wing). The partial
differential equations for viscous fluid flow are quite complicated, but he argued that under certain conditions the effects of
viscosity are concentrated in a narrow layer near the surface of the body. This happens, for example, with air flow across an
airplane wing. This observation allowed Prandtl to go through an order of magnitude argument and omit terms he felt to
be negligible in the equations. The result was a problem that he was able to solve. This was a brilliant piece of work, but
it relied strongly on his physical intuition. For this reason, there were numerous questions about his reduction that went
unresolved for decades. For example, it was unclear how to obtain the correction to his approximation, and it is now thought
that Prandtl’s derivation of the second term is incorrect (Lagerstrom, 1988). This predicament was resolved when Friedrichs
(1942) was able to show how to reduce a boundary layer problem systematically. In analyzing a model problem, he used
a stretching transformation to match inner and outer solutions, which is the basis of the method discussed in this chapter.
This procedure was not new, however, as demonstrated by Gans’ (1915) use of some of these ideas to solve problems in
optics.
The golden age for matched asymptotic expansions was the 1950s, and it was during this period that the method was
refined and applied to a wide variety of physical problems. A short historical development of the method is presented in
O’Malley (1991) and a nice discussion of boundary layer methods used in the nineteenth century is given by Van Dyke (1994).
The popularity of matched asymptotic expansions was also greatly enhanced with the appearance of two very good books,
one by Van Dyke (1964) and another by Cole (1968). The method is now one of the cornerstones of applied mathematics.
At the same time, it is still being extended, both in the type of problems it is used to resolve as well as in the theory.
Generally speaking, the matched asymptotics technique tends to be most useful in the solution of certain boundary
value problems though some initial value problems are also amenable to this approach.
5.2 Rescaling differential equations: general considerations

We have already discussed rescaling non differential equations to bring out different balances. When dealing with differen-
tial equations, similar considerations apply but the procedure is more complicated because there are both dependent and
independent variables e.g. y(x). Much of the present section will be concerned with rescaling the independent variable, and
in particular we will introduce stretching transformations of the independent variable x. For linear homogeneous problems
42
(equation and BCS), we usually do not need to rescale the dependent variable y for the simple reason that any rescaling will
have no effect on the overall balances in the equation. For example consider the equation:
y ′′ + y ′ + y = 0. (153)
If we rescale via y(x) = εY (x), it is clear that the equation for Y (x) will reflect exactly the same balances; however while
the balances are not changed by rescaling, it is possible that the rescaled variable is suggested by a boundary condition or
during asymptotic matching. If however the equation is non-linear, then rescaling the dependent variable can change the
balances. For example consider:
y ′′ + yy ′ + y = 0. (154)
If we rescale via y(x) = εY (x), then (on dividing through by ε) the equation becomes:
Y ′′ (x) + εY (x)Y ′ (x) + Y (x) = 0. (155)

and this does reflect a new balance. Similar considerations apply when dealing with inhomogeneous linear equations. It
should be borne in mind that boundary conditions must also be considered when one is rescaling; when one is rescaling, one
is considering the whole problem and a new balance may be reflected in the boundary conditions but not in the differential
equation. For example, consider the problem:
y ′′ + y ′ + y = 0; y(0)2 − y ′ (0) = 0, y(1) = 0. (156)

The rescaling y(x) = εY (x) does not change the balances in the differential equation but it does change the boundary
conditions (because they contain a non-linear term). Thus the rescaled problem is:
Y ′′ + Y ′ + Y = 0; εY (0)2 − Y ′ (0) = 0, Y (1) = 0. (157)

Matters become even more complicated when we deal with systems of differential equations and/or partial differential
equations. The question of how to determine the which particular rescaling is appropriate is a key issue in solving differential
equations with a small parameter.
In summary, if a boundary value problem (including its boundary conditions) is nonlinear or inhomogeneous, it may sometimes
be necessary to rescale the dependent variable y.
5.3 Example
The best way to explain the method of matched asymptotic expansions is to use it to solve a problem. The following example
takes several pages to complete because it is used to introduce the ideas and terminology. As the procedure becomes more
routine, the derivations will become much shorter.
The problem we will study is the boundary value problem:
εy ′′ + 2y ′ + 2y = 0, for 0 < x < 1, (158)

where
y(0) = 0, y(1) = 1 (159)

This equation is similar to the one used in the previous section to discuss uniform and nonuniform approximations. The
difference is that we will now derive the approximation directly from the problem rather than from a formula for the solution.
An indication that this problem is not going to be as straightforward as the differential equations considered previously is
that if ε = 0 then the problem is no longer second order. This leads to what is generally known as a singular perturbation
problem, although singularity can occur for other reasons. In any case, to construct a first-term approximation of the
solution for small ε we will proceed in four steps. The fifth step will be concerned with the derivation of the second term in
the expansion.
5.3.1 Step 1: Outer Solution

To begin, we will assume that the solution can be expanded in powers of ε. In other words,
y(x) ∼ y0 (x) + εy1 (x) + . . . . (160)

Substituting this into (158), we obtain
ε(y0′′ + . . .) + 2(y0′ + εy1′ + . . .) + 2(y0 + εy1 + . . .) = 0 (161)

with boundary conditions
y0 (0) + εy1 (0) + ... = 0; y0 (1) + εy1 (1) + ... = 1. (162)
43
ε0 problem The ε0 problem (collecting terms ε0 in (161)) is therefore
y0′ + y0 = 0, y0 (0) = 0; y0 (1) = 1 (163)
ε1 problem The ε1 problem (collecting terms ε1 in (161))is :
y1′ + y1 = −y0′′ , y1 (0) = 0; y1 (1) = 0 (164)

1
Note that the ODE at ε is inhomogeneous and dealing with the right hand side involves first solving the zero order
equation. Note also that the boundary conditions at ε1 are different than those at ε0 . In the present instance we will pay
most attention to the leading order solution in order to demonstrate the general technique. However we return to the ε1
problem later.
It is useful to note that the ε0 problem is just the original problem with ε = 0. The general solution is
y0 (x) = ae−x , (165)

where a is an arbitrary constant. Looking at the solution in (165), we have a dilemma because there is only one arbitrary
constant but there are two boundary conditions (159). What this means is that the solution in (165) and the expansion in
(160) are incapable of describing the solution over the entire interval 0 ≤ x ≤ 1. At the moment, we have no idea which
boundary condition, if any, we should require y0 (x) to satisfy, and the determination of this will have to come later. This
leads to the question of what to do next. Well, based on what is observed in Example 2 in Section 1.7, it is a reasonable
working hypothesis to assume that (165) describes the solution over most of the interval but that there is a boundary layer
at either x = 0, or at x = 1, where a different approximation must be used. Assuming for the moment that it’s at x = 0,
then we are probably missing a term like e−βx/ε . The derivation of this term is the objective of the next step. Because we
are going to end up with approximations of the solution over different regions, we will refer to (165) as the first term in the
expansion of the outer solution.
5.3.2 Step 2: Rescaling the problem: the boundary layer

Based on the assumption that there is a boundary layer at x = 0, we anticipate that the ODE behaves in a different way
near x = 0 in that different terms form the dominant balance in that region. In order to bring out the new balances in
the equation we rescale the problem. (In general one scales to bring our particular balances in problems, both in the
equations and in the boundary conditions) Thus we assume that the original non-dimensionalisation of the problem which
led to (158) does not represent a good scaling of the problem near x = 0. In particular we expect that the second derivative
term of (158) may become significant near x = 0 and this provides some guidance as to how to rescale. (As the equation
and boundary conditions are linear and homogeneous, we will certainly not have to rescale the dependent variable y). In the
present case it is thus sufficient to rescale the independent variable so we introduce a boundary-layer coordinate given as
x = 0 + εα x. (166)
In this problem, the boundary layer is located near x = 0 but in other problems it may be located elsewhere and it is useful
to get into the habit of inserting the location explicitly. Thus
x
x= , (Y (x) = y(x)) (167)
εα
where α > 0. From our earlier discussion, it is expected that α = 1, and this will be shown conclusively below via balancing
of terms. After changing variables from x to x, we will take x to be fixed when expanding the solution in terms of ε. This
has the effect of stretching the region near x = 0 as ε becomes small. Because of this, (167) is sometimes referred to as a
stretching transformation. From the change of variables in (167), and the chain rule, we have that
d dx d 1 d
= = α .
dx dx dx ε dx
If we let Y (x) denote the solution of the problem when using this boundary-layer coordinate, then (158) transforms to
d2 Y −α dY
ε1−2α 2 + 2ε + 2Y = 0, (168)
dx dx
where, from (159),
Y (0) = 0. (169)
(168) and (169) comprise the rescaled problem, which is only valid near x = 0. Thus the other boundary condition will
not in general be appropriate. We will show below that α = 1 and we comment that (168) demonstrates a fundamentally
different behaviour from (158) in terms of asymptotic balances.
44
The boundary condition at x = 0 has been included here because the boundary layer is at the left end of the interval.
The appropriate expansion for the boundary-layer solution is now
Y (x) ∼ Y0 (x) + ε1 Y1 (x) + . . . , (170)

As stated above, in this expansion x is held fixed as ε goes to zero (in the same way that x is held fixed in the outer
expansion). Substituting the expansion in (170) into (168), we get that
d2 d
ε1−2α (Y0 + . . .) + 2ε−α (Y0 + . . .) + 2(Y0 + . . .) = 0. (171)
dx2 dx
1 2 3
Just as with the algebraic equations studied in Section 2.4, it is now necessary to determine the correct balances in (171).
The balance between terms 2 and 3 was considered in Step 1, and so the following possibilities remain (Exercise):
i) 1 ∼ 3 and 2 is higher order ⇒ 1 − 2α = 0 ⇒ α = 21 ⇒ 1 , 3 = O(1) but 2 = O(ε−1/2 ). This violates our original
assumption that 2 is higher order, and so this balance is not possible.
ii) 1 ∼ 2 and 3 is higher order ⇒ 1 − 2α = −α ⇒ α = 1 ⇒ 1 , 2 = O(ε−1 ) and 3 = O(1). In this case, the
conclusions are consistent with the original assumptions, and so this is the balancing we are looking for. This is said to be
a distinguished limit for the equation. Thus
The rescaled problem In terms of the new variables Y (x) we have
d2 d
2 (Y0 + εY1 . . .) + 2 dx (Y0 + εY1 . . .) + 2ε(Y0 + εY1 . . .) = 0; Y0 (x = 0) + εY1 (x = 0) = 0. (172)
dx
This leads to the following sequence of problems:
ε0 problem
ε0 Y0′′ + 2Y0′ = 0, for 0 < x < ∞, (173)
Y0 (x = 0) = 0. (174)
ε1 problem
ε1 Y1′′ + 2Y1′ = −2Y0 , for 0 < x < ∞, (175)
Y1 (x = 0) = 0. (176)
Notice that the ε1 ODE is inhomogeneous. Solution of it requires the ε0 solution.
The general solution of the ε0 (linear, second order, constant coefficient ODE) problem is
Y0 (x) = A(1 − e−2x ), (177)

where A is an arbitrary constant. It should be observed that the boundary-layer equation (173) contains at least one term
of the outer-layer equation (163). This is important for the successful completion of Step 3.
The boundary layer expansion in (170) is supposed to describe the solution in the immediate vicinity of the endpoint
x = 0. It is therefore not unreasonable to expect that the outer solution (165) applies over the remainder of the interval (this
is assuming there are no other layers). This means that the outer solution should satisfy the boundary condition at x = 1.
From (165) and (159), one finds that
y0 (x) = e1−x . (178)
Step 3: Matching It remains to determine the constant A in the first-term approximation of the boundary-layer solution
(177). To do this, the approximations we have constructed so far are summarized in Fig. 16. The important point here
is that both the inner and outer expansions are approximations of the same function. Therefore, in the transition region
between the inner and outer layers we should expect that the two expansions give the same result. This is accomplished by
requiring that the value of Y0 as one comes out of the boundary layer (i.e., as x → ∞) is equal to the value of y0 as one
comes into the boundary layer (i.e., as x → 0). Imposing this condition yields Y0 (∞) = y0 (0).
Another way of thinking of the matching: both the inner and outer expansions are valid in the (overlap) region ε ≪ x ≪ 1,
corresponding to x → 0 and x → ∞ as ε → 0.
They satisfy the matching condition Thus, A = e1 and (177) becomes
45
A
e1
Solution
Outer Approximation
Inner Approximation
10-5 10-4 10-3 10-2 10-1 1

x - axis
Figure 16: Sketch of the inner solution (177) and the outer solution (178).Note the overlap region along the x− axis where
both solutions are essentially constant. Since these approximations are supposed to be describing the same continuous
function, it must be that these constants are the same. So A = e.
Y0 (x) = e1 − e1−2x . (179)

This completes the derivation of the inner and outer approximations of the solution of (158). The next step is to combine
them into a single expression (this is Step 4). Note that this approach to matching will work when we are working with one
term expansions. For higher order approximations it is convenient to introduce alternative matching schemes.
Important note An important practical point is that the inner problem is effectively solved on an infinite domain.
Suppose that the outer problem has a boundary condition specified at x = 1 as in this problem where y(1) = 1. In terms
of the inner variables, this translates into a condition at the far end x = 1/ε. However, this equals infinity for all practical
purposes, and the condition at x = 1/ε is replaced by a matching condition at infinity in effect (in the inner region). This
feature is common to all boundary layer problems.
Matching in more detail
Van Dyke Matching This procedure is relatively simple to use but can occasionally lead to incorrect results (Fraenkel,
1969; Lagerstrom, 1988).
This method uses the following procedure: suppose the outer solution y outer (x) or y o (x), with outer variable x, is obtained
to O(∆(ε)) and the inner solution y inner (x) or y i (x), with inner variable x, is known to O(δ(ε)).
The outer solution is written in terms of the inner variable x which is then expanded in terms of ε ≪ 1 to O(δ) with all
variables assumed O(1) and then truncated. This is called the inner expansion of the outer solution (y o )i . The inner solution
is written in terms of the outer variable x, which is then expanded to O(∆) and then truncated. This is called the outer
expansion of the inner solution (y i )o . Matching requires (y o )i = (y i )o .
Note finally that if the dependent variable has been rescaled, it is important when doing the matching that this is
taken into account i.e., for the purposes of the matching both the inner and outer solutions should be written in terms of
either the inner or outer dependent variable.
The matching idea The idea is that different expansions for the same function almost always coincide in overlap regions.
Although the van Dyke matching principle is not universal, it has the advantage of being easy to apply. If it fails, one must
use intermediate variables. To get some feeling for the principle, consider the function 1/(x + ε) and expand this for small ε
with x = O(1). (Consider this an outer solution). The expansion fails near x = 0 where the rescaling x = εx̄ is appropriate
whence we deal with 1/ε(x̄ + 1) with x̄ = O(1). (Consider this an inner solution). Now try applying the matching principle
to the above “inner” and “outer” solutions.
46
The outer solution expanded to O(1) is:
1 1
∼ + O(ε). (180)
x+ε x
The inner solution is 1/ε(x̄ + 1) and its expansion with x̄ held constant is just itself i.e.,
1 1
. (181)
ε (x̄ + 1)
We thus regard this as an inner expansion being known to O(1/ε). Now applying van Dyke we need to expand the outer
solution as far as O(1/ε) i.e.,
1 1
= . (182)
x εx̄
Similarly we expand the inner solution in terms of the outer variable as far as O(1). Thus:
1 1 x 1
= ε−1 ( + 1)−1 = (x + ε)−1 ∼ + O(ε). (183)
ε (x̄ + 1) ε x
Matching now requires that

1 1
= (184)
εx̄ x
which is indeed the case as x = εx̄.
In words, the to O(δ) inner expansion of the to O(∆) outer solution equals the to O(∆) outer expansion of the to O(δ) inner
solution.
Our outer solution from (178) is y o = e1−x and as we only know one term of it, ∆ = 1. Our inner solution from (177) is
y i = A(1 − e−2x ) and so δ = 1. To obtain (y o )i , we subsitute x = εx into y o and (based on ε ≪ 1) expand to O(1) but no
further. Thus
e1−x = e1−εx = e − εx + O(ε2 ) ∼ e.
To obtain (y i )o , we substitute x = x/ε into y i and expand to O(1).Thus
A(1 − e−2x ) = A(1 − e−2x/ε ) = A + EST ∼ A

where EST stands for exponentially small terms i.e., terms which are smaller than any power of ε.(Recall that e−1/ε ≪ εn
as ε ↓ 0). Thus for matching to occur we require that:
A=e
and the unknown constant in the inner solution y i = e(1 − e−2x ) has been determined. Note that (y o )i and (y i )o are referred
to as being the common part of the inner and outer solutions. They represent the behaviour of both the inner and outer
solutions in an intermediate region.
Intermediate variable matching An alternative to van Dyke matching is via intermediate variables. The approach
to be described will work on more difficult problems than the simple example in (158). The idea is to introduce an immediate
variable xη = x/η(ε) that is positioned between the O(1) coordinate of the outer layer and the O(ε) coordinate of the inner
layer. This intermediate variable is to be located in the transition region, or overlap domain, seen in Fig. 16, and for this
reason it should be expected that η(ε) satisfies ε ≪ η ≪ 1. The precise conditions imposed on η(ε) are stated explicitly in
the following matching procedure (see Fig. 17):
i) Change variables in the outer expansion (from x to xη ) to obtain youter . It is assumed that there is an η1 (ε) so that
youter still provides a first-term expansion of the solution for any η(ε) that satisfies η1 (ε) ≪ η(ε) ≤ 1.
ii) Change variables in the inner expansion (from x to xη ) to obtain yinner . It is assumed that there is an η2 (ε) so that
yinner still provides a first-term expansionof the solution for any η(ε) that satisfies ε ≤ η(ε) ≪ η2 .
iii) It is assumed that the domains of validity of the expansions for youter and yinner overlap, that is, η1 ≪ η2 . In this
overlap region, for the expansions to match, it is required that the first terms from youter and yinner are equal.
The assumptions contained in (i) and (ii) can be proved under fairly mild conditions and are the essence of Kaplun’s
extension theorem (Lagerstrom, 1988). Actually, one seldom is interested in determining n1 or n2 but only that there is an
interval for η(ε) so that yinner and youter match. It is important, however, that the matching not depend on the specific
choice of η(ε). For example, if one finds that matching can only occur if η(ε) = ε1/2 , then there is no overlap domain and
the procedure is not applicable. In comparison to the situation for (i) and (ii), the assumption on the existence of an overlap
47
Outer
Inner
0 ε η1 η η2 1 x - axis
Overlap Domain
Figure 17: Schematic of the domains of validity of the inner and outer expansions as assumed in the matching procedure.
The intermediate variable is to be located within the overlap region.
domain in (iii) is a different matter, and a satisfactory proof has never been given. For this reason, it has become known as
Kaplun’s hypothesis on the domain of validity (Lagerstrom, 1988).
To use the matching procedure just described, we introduce the intermediate variable xη , defined as
x
xη = , (185)
εβ
where 0 < β < 1. This interval for β comes from the requirement that the scaling for the intermediate variable lies between
the outer scale, O(1), and the inner scale O(ε). Actually, it may be that in carrying out the matching of yinner and youter
we have to reduce this interval for β. In any case, to see if the expansions match, note that the inner solution, from (170)
and (177), becomes
1−β
yinner ∼ A(1 − e−2xη /ε ) + ... (186)
∼ A+ ...,
and the outer solution, from (160) and (178), becomes
β
youter ∼ e1−xη ε + . . . (187)
1
∼ e + ...,
The expansions in (186) and (187) are supposed to agree to the first term in the overlap domain and therefore A = e1 .
The modification of the intermediate variable procedure to deal with higher order expansions will not be given here, but
in Step 5 a two-term expansion is matched to illustrate the situation. For those who may be interested, Lagerstrom (1988)
has a more extensive discussion of the subtleties of matching using an intermediate variable.
5.3.3 Step 4: Composite expansion

Our description of the solution consists of two pieces, which we will now combine to form a composite expansion. This is
done by adding the expansions and then subtracting the part that is common to both. Thus,
x
y ∼ y o (x) + y i − common part (188)
ε
∼ e1−x − e1−2x/ε .
It may seem strange that it is possible to combine expansions over different intervals and still have an asymptotic approx-
imation of the solution. However, note that in the outer region the boundary-layer solution, Y0 (x/ε), is constant to first
order. This constant is y0 (0), and to compensate for its contribution the term −y0 (0) is included in the composite expansion.
Similarly, in the boundary-layer region the outer solution, y0 (x), is constant to first order. However, the −y0 (0) term removes
its contribution in this region. The fact that the adjustment in each case involves the constant y0 (0) is not a coincidence
since it is the first term in the inner and outer expansions found from matching. Because of this, it is referred to as the
common part of the expansions.
48
3.0
2.0
Solution
1.0 Exact Solution

Composite Expansion
0 0.2 0.4 0.6 0.8 1.0

x - axis
Figure 18: Graph of the exact solution of (158) and the composite expansion given in (188) with ε = 0.01. The two curves
are almost indistinguishable.
We can regard the outer expansion of the inner solution, y io , to be the important part of y i in the outer region. Thus
y − y io should be negligible in the outer region. Similarly yo − y oi should be negligible in the inner region. Thus, using these
i
interpretations, we can write the composite solution as
yc = y o + (y i − y io ) = y i + (y o − y 0i ).
The fact that the composite expansion gives a very good approximation to the solution over the entire interval is shown
in Fig 18. Note, however, that it satisfies the boundary condition at x = 0 exactly but the one at x = 1 is only satisfied
asymptotically. This is not of particular concern since the expansion also satisfies the differential equation in an asymptotic
sense. However, an alternative expansion, which satisfies both boundary conditions, is developed in the exercises.
5.3.4 Step 5: Second term

Generally, to illustrate a method we will only derive the first term in an expansion. However, the second term is important
as it gives a measure of the error. The procedure to find the second term is very similar to finding the first so only the
highlights will be given here.
By substituting the outer expansion (160) into the problem and collecting the ε1 terms, one finds that y1′ + y1 = − 21 y0′′
with y1 (1) = 0. as in (164). The solution of this problem is
1
y1 = (1 − x) e1−x .
2
Similarly, from the boundary-layer equation at the ε1 order (175) we get that Y1′′ + 2Y1′ = −2Y0 with Y1 (0) = 0. The general
solution of this is

Y1 = B 1 − e−2x − xe1 1 + e−2x (189)
where B is an arbitrary constant.
5.3.5 Higher order terms

van Dyke 2 term matching In this case we have the outer and inner expansions both to O(ε) so δ = ∆ = ǫ. The outer
solution is
ε
youter (x) = e1−x + (1 − x)e1−x (190)
2
while the inner solution is:
49

yinner = A(1 − e−2x ) + εB 1 − e−2x − εxe1 1 + e−2x . (191)
In fact,we have already shown that A = e from the 1 term matching. We show now that the same result obtains during
2 term matching. We write (190) in terms of the inner variable with x = εx and expand for ε ≪ 1 as far as O(ε) assuming
that x = O(1) Thus
ε ε
(youter )i = e1−ex + (1 − εx)e1−εx = e1 (1 − εx + ..) + e1 + O(ε2 ) (192)
2 2
while writing (191) in terms of the outer variable with x = x/ε and expanding for ε ≪ 1 as far as O(ε) assuming that
x = O(1) yields:
x x x
(yinner )o = A(1 − e−2 ε ) + εB 1 − e−2 ε − xe1 1 + e−2 ε = A + εB − xe + O(ε2 ) (193)
To compare (192) and (193), we can either write the latter in terms of x or the former in terms of x. Substituting x = εx
into (193), we get:
A + εB − εxe (194)
and comparing this to (192), we see that we must choose:
A = e, B = e/2. (195)
and the other term matches automatically. (It is important to check that this is so). To this order, the common part of the
two expansions is (194): this needs to be subtracted off to form the composite solution.
Occasionally it happens that an expansion (inner or outer) produces a term of an order, or form, that the other does not
have. A typical example of this occurs when trying to expand in powers of ε. It can happen that to be able to match with an
expansion from an adjacent layer it is necessary to include other terms, such as those involving ln(ε). This process of having
to insert scales into an expansion because of what is happening is another layer is called switchbacking. Some of the more
famous examples of this involve logarithmic scales, and these are discussed in Lagerstrom (1988).
Intermediate variable 2 term matching To match the expansions, we use the intermediate variable given in (185).
The outer expansion in this case takes the form
εβ β
youter ∼ e1−xη ε + 1 − xη εβ e1−xη ε + . . . (196)
2
ε 1
1 β 1
∼ e − ε xη e + e1 + ε2β e1 x2η + . . . ,
2 2
and, setting ξ = −2xη /ε1−β , the boundary-layer expansion becomes

1 ξ
ξ
xη e1 ξ

yinner ∼ e 1−e + ε B 1 − e − 1−β 1 + e + ... (197)
ε
∼ e1 − εβ xη e1 + Bε.
Matching these we get B = 21 e1 . Note, however, that these expansions do not appear to agree since (196) contains a
O(ε2β ) term that (197) does not have. To explain why this occurs, note that both expansions produce a O(εβ ) term that
does not contain an arbitrary constant. If this term is not identical for both expansions, there is no way the expansions will
match. In the outer expansion, this term comes from the O(1) problem, and in the boundary layer it comes from the O(ε)
solution. In a similar manner, one finds that the x2η term in (196) also comes from the first term. However, for the boundary
layer it comes from the O(ε2 ) problem (the verification of this is left as an exercise). Therefore, the expansions match.
Higher order composite expansion It is now possible to construct a two-term composite expansion. The basic idea is
to add expansions and then subtract the common part. In the case of van Dyke matching this yields the result:
e
y ∼ y0 + εy1 + Y0 + εY1 − e + ε − εxe
2
1−x 1−2x/ε εh i
∼e − (1 + x) e + (1 − x) e1−x − e1−2x/ε .
2
This solution is valid everywhere in the domain with an O(ε2 ) error.
If we use intermediate variable matching, we obtain the same result:
50
√ ε
y ∼ y0 + εy1 + Y0 + εY1 − e1 − xη e1 ε + e1
2
1−x 1−2x/ε εh 1−x
i
∼e − (1 + x) e + (1 − x) e − e1−2x/ε .
2
Note that the common part in this case contains the terms in (196) and (197) except for the x2η term in (196).
5.3.6 Discussion
The importance of matching cannot be overemphasized. Numerous assumptions went into the derivation of the inner and
outer approximations, and matching is one of the essential steps that supports these assumptions. If they had not matched,
it would have been necessary to go back and determine where the error was made. The possibilities when this happens are
almost endless, but the following are useful places to start to look:
1. The boundary layer is at x = 1 and not at x = 0. In this case, the boundary-layer coordinate is
x−1
x= or x = 1 + εα x. (198)
εα
For certain problems, it may be necessary to replace the denominator with a function µ(ε), where µ(ε) is determined
from the balancing or matching.
2. There are boundary layers at both ends of the interval. See later.
3. There is an interior layer. In this case, the stretching transformation is
x − x0
x= or x = x0 + εα x. (199)
εα
where x0 is the location of the layer (it may depend on ε). See later.
4. The form of the expansion is incorrect. For example, the outer expansion may have the form y ∼ ξ(ε)y0 (x) + . . . , where
ξ is a function determined from the balancing or matching in the problem.
5. The solution simply does not have a layer structure, and other methods need to be used.
Occasionally, it is so unclear how to proceed that one may want to try to solve the problem numerically to get an insight into
the structure of the solution. The difficulty with this is that the presence of boundary or interior layers can make it hard
to obtain an accurate numerical solution (see Section 2.7). If this happens, one can try taking values of ε that are not too
small so the layer is present but not so localized that it dominates the numerical solution. Another possibility to help guide
the analysis occurs when the problem originates from an application and one is able to use physical intuition to determine
the approach to take to solve the problem.
5.4 Example
Consider the following problem for a function u(y):
d2 u du
ǫ + = 1; u(0) = 0, u(1) = 2, (200)
dy 2 dy
where ǫ denotes a small positive constant. Note that this is a linear inhomogeneous ODE with constant coefficients and
we can solve it exactly. The exact solution can be written in the form u = uc + up where uc is the solution of the related
homogeneous problem ǫu′′ + u′ = 0 and up is any particular solution of the full problem. Noting that up = y is a particular
solution, the exact solution (fig. (19)) is easily (?) shown to be
1 − e−y/ǫ
u=y+ . (201)
1 − e−1/ǫ
Suppose however that we were unable to solve (200) exactly. In view of the small parameter ǫ ≪ 1, is it possible to obtain
an approximate solution? Let us examine this while bearing in mind that if we get stuck we can always examine the exact
solution to see what is going on. If this were a real problem which we could not solve exactly, we would not have this luxury,
of course.
Now, e−1/ǫ is extremely small, and so is e−y/ǫ , for 0 < y < 1, unless y = O(ǫ). The solution may therefore be
approximated, in two parts, by a ‘mainstream’
51
uM = y + 1, (202)
and a ‘boundary layer’ adjacent to y = 0 with thickness O(ǫ):
uBL = 1 − e−y/ǫ . (203)

These two expressions represent particular limits of the full, exact solution (201), the first being obtained by letting ǫ → 0
for fixed y, and the second being obtained by letting ǫ → 0 for y/ǫ fixed. Let us now be more precise about this and use the
method of matched asymptotic expansions.
5.4.1 Approximate solution using the method of matched asymptotic expansions

We first look for a straightforward 2 term expansion of (200) by letting u = u0 (y) + εu1 (y) and we find that
u′0 = 1 =⇒ u0 = y + c. (204)
The relevant boundary conditions at y = 0 become:
u0 (0) + εu1 (0) = 1 =⇒ u0 (0) = 1, u1 (0) = 0 (205)

and so we find that:
u0 (y) = y + 1. (206)
At the next order we find that u1 (y) must satisfy:
′
u1 = 0 (207)
and from the boundary condition we see that u1 (y) = 0. Thus a 2 term outer expansion is
u = y + 1 + O(ε2 ) (208)
α
In the inner region, we introduce a stretched inner coordinate via Y = y/ε and we quickly find that α = 1 by examining
the relevant balances in the equation. Hence the inner independent variable is:
Y = y/ǫ. (209)
With this scaling the previously negligible second derivative regains its importance and the inner equation is:
1 d2 U 1 dU
ǫ 2 2
+ = 1. (210)
ǫ dY ǫ dY
The inner boundary condition is:
U (0) + εU1 (0) = 0 =⇒ U (0) = 0, U1 (0) = 0. (211)

Thus we find that
d2 U0 dU0
+ = 0. (212)
dY 2 dY
This is a second order, linear, constant coefficient ODE. It can be solved using the exponential substitution U0 = eλ y to find
two linearly independent solutions, or first integrated once to reduce to a first order linear ODE which can be solved directly.
On making the inner solution satisfy the boundary condition U0 (0) = 0 we obtain
U0 = A(1 − e−Y ), (213)

where A is a constant to be determined in the matching. Proceeding to the next order we find that:
d2 U1 dU1
+ = 1. (214)
dY 2 dY
Noting that U1 = Y is a particular solution of this equation we find the general solution to take the form :
U1 = E + F e−Y + Y (215)
and applying the BC U1 (0) = 0 we find that
U1 = E(1 − e−Y ) + Y (216)

2
so the inner solution to O(ε ) is:
U = A(1 − e−Y ) + ε[E(1 − e−Y ) + Y ] (217)
52
2.0
1.8
1.6
1.4
outer
1.2 composite
u
1.0
0.8
inner
0.6
0.4
0.2
0.0
0.0 0.1 0.2 0.3 0.4 0.5 0.6 0.7 0.8 0.9 1.0
Figure 19: Inner, outer and composite solutions for ε = 0.1. The exact solution is indistinguishable from the composite
solution.
5.4.2 Matching
Using van Dyke matching, we note that we have an inner and outer solution both to O(ε), so in this case δ = ∆ = ε.Thus
taking (208) in terms of the inner variable, we obtain
(u)i = εY + 1 (218)
and expanding this for ε ≪ 1, Y = O(1) to O(ε) and truncating leaves it unchanged. Writing the inner solution (217) in
terms of the outer variable, we obtain
(U )o = A(1 − e−y/ε ) + ε[E(1 − e−y/ε ) + y/ε] (219)

Now expanding this for ε ≪ 1, y = O(1) we note the presence of E.S.T. which we neglect and we obtain:
(U )o = A + y + εE (220)
To compare (218) and (220), we must use the same variable so we set y = εY in (220) and we see that we require that
(U )o = (u)i . Hence
εY + 1 = A + εY + εE (221)
so we must choose A = 1, E = 0 to complete the matching. Finally we write a composite series as u + U − (U )0 where
(U )o = 1 + y
u = (y + 1) + ((1 − e−y/ε ) + y) − (1 + y) = y + 1 − e−y/ε . (222)

Note that we chose to write the solution in terms of the outer variable y; we could as easily have used the inner variable Y.
Figure 19 shows the inner, outer and composite solutions.
The basic idea behind matched asymptotics Suppose we have a differential equation problem (equation + boundary
conditions) which contains a small parameter (ε ≪ 1) and we know the exact solution. Then we expect that we can obtain
an approximate solution by exploiting the smallness of ε and expanding in an appropriate manner.
If we do not know the exact solution, perhaps we can obtain an approximate solution, by expanding the equations and
boundary conditions in an appropriate way and solving the associated simpler problems. The equation+BCs must implicitly
contain all the information about the unknown solution so this approach sounds at least plausible.
For example (200) has exact solution (201). However we cannot get an approximate version of (201) in one step. One
case is to take the case y = O(1) and consider the limit ε → 0 which leads to the approximate u ∼ y + 1 which is (202).
This limiting process fails when y = O(ε) so that exp (−y/ε) is no longer small. The other case thus deals with the situation
where y = O(ε). In this case we let Y ≡ y/ε = O(1) and so consider the limit ε → 0 with y/ε held fixed (i.e., O(1)). This
53
leads to the the solution U = 1 − exp (−y/ε) which is (203). So the exact solution has two distinct behaviours: y fixed, ε → 0
and y/ε fixed, ε → 0.
Now suppose that we do not have the exact solution. Then our approach is to consider the whole problem in the two
appropriate limits y fixed, ε → 0 and y/ε fixed, ε → 0. (These are the limits one considers when one rescales using the
variables y and y/ε. This gives rise to two approximate sub-problems: the inner and outer problems: this is how the method
of matched asymptotics works. It tries to identify what the appropriate limiting behaviours are for the exact solution and
applies these limits to the equation and BCs rather than to the exact solution. The appropriate limiting bahaviours are the
scales which we use in the inner and outer regions.
5.5 Multiple boundary layers

Almost all of the principal ideas underlying the method of matched asymptotic expansions were introduced in the example
of section 5.3. In the remainder of this chapter, these ideas are applied, and in most cases extended, to more complicated
problems. The example we now consider illustrates the situation of multiple boundary layers. The problem is
ε2 y ′′ + εxy ′ − y = −ex , for 0 < x < 1, (223)

where
y(0) = 2 and y(1) = 1. (224)

One difference here is that one of the coefficients of the equation depends on x. This is not responsible for the multiple
boundary layers but, as we shall see, it does result in different equations for each layer.
5.5.1 Step 1: Outer Expansion

The expansion of the solution in this region is the same as in the last example, namely, y ∼ y0 + . . . . From this and (223)
one obtains the first-term approximation
y0 = ex . (225)
Clearly, this function is incapable of satisfying either boundary condition, an indication that there are boundary layers at
each end.
5.5.2 Steps 2 and 3: Boundary Layers and Matching

For the left endpoint we introduce the boundary-layer coordinate x = 0 + εα x or equivalently x = x/εα , in which case (223)
becomes
d2 Y dY α
ε2−2α + εx − Y = −eε x . (226)
dx2 dx
1 2 3 4
As before, Y (x) is used to designate the solution in this boundary-layer region. The balance in this layer is between terms
1 , 3 , 4 , and so α = 1. The appropriate expansion for Y is Y ∼ Y0 (x) + . . . , and from (226) and (224) we have that
Y0′′ − Y0 = −1, for 0 < x < ∞, (227)

where
Y0 (0) = 2. (228)
Note that the ODE is linear, second order, constant coefficient, inhomogeneous. Note also that (227) has at least one term
in common with the equation for the outer region (which we should expect if there is to be any hope of matching the inner
and outer expansions). The general solution of (227) is
Y0 (x) = 1 + Ae−x + (1 − A)ex . (229)

This must match with the outer solution given in (225). The matching condition is Y0 (∞) = y0 (0), and so from (225) we have
that A = 1. To determine the solution in the boundary layer at the other end, we introduce the boundary-layer coordinate
x−1
x
e= ; (x = 1 + εβ x
e). (230)
εβ
In this region, we will designate the solution as Ye (e
x). Introducing (230) into (223), one obtains the equation
54
3.0
Numerical Solution
2.5
Composite Expansion
Solution
2.0
1.5
1.0
0.2 0.4 0.6 0.8 1.0
x - axis
Figure 20: Composite expansion (235) and the numerical solution. ε = 0.01. The curves are almost indistinguishable.
d2 Ye dYe β
ε2−2β 2
x)ε1−β
+ (1 + εe − Ye = −e1+ε xe. (231)
de
x dex
The relevant distinguished limit in this case occurs when β = 1. So, the expansion Ye ∼ Ye0 (e
x) yields the problem
Ye0′′ + Ye0′ − Ye0 = −e1 , for − ∞ < x

e < 0, (232)
where
Ye0 (0) = 1. (233)

It is important to notice that this boundary-layer equation has at least one term in common with the equation for the outer
region. Solving (232), noting that it is a linear, second order, constant coefficient, inhomogeneous ODE, the general solution
is found to be
+ −
Ye0 (e
x) = e1 + Ber xe + (1 − e1 − B)er xe, (234)
√
where 2r± = −1 ± 5. The matching requirement here is that Ye0 (−∞) = y0 (1). (Exercise: do the matching using van Dyke’s
technique noting for example that r+ > 0, r− < 0 when x < 1). Hence, from (234), B = 1 − e1 . As a final comment, note that
boundary-layer equation (232) differs from (227), and this is due to the x dependence of the terms in the original problem.
5.5.3 Step 4: Composite Expansion

The last step is to combine the three expansions into a single expression. This is done in the usual way of adding the
expansions together and then subtracting the common parts. From (225), (229), and (234), a first term expansion of the
solution over the entire interval is
y ∼ y0 (x) + Y0 (x) − Y0 (∞) + Ye0 (e

x) − Ye0 (−∞) (235)
x −x/ε 1 −r(1−x)/ε
∼e +e + (1 − e )e ,
√
where 2r = −1 + 5. This approximation is shown in Fig.20 along with the numerical solution. One is able to clearly see
the boundary layers at the endpoints of the interval as well as the accuracy of the asymptotic approximation.
55
5.5.4 Comments
Some interesting complications arise when using matched asymptotic expansions with nonlinear equations. For example, it is
not unusual that the solution of a nonlinear equation is defined implicitly. To illustrate this situation, consider the problem
εy ′′ + εy ′ − ey = −2 − x, for 0 < x < 1, (236)

where y(0) = 0 and y(1) = 1. The
√ outer expansion is y ∼ ln(x + 2), and this holds for 0 < x < 1. For the boundary layer at
x = 0 the coordinate is x = x/ ε, and one finds that Y ∼ Y0 , where
Y0′′ − eY0 = −2,

for Y0 (0) = 0. Note that this equation is autonomous i.e., x does not appear explicitly. Thus one can reduce the order of
the equation by one by using the substitution1
′ du0 dY0 ′
u(Y0 ) = Y0 =⇒ Y0′′ = = u0 (Y0 )u0 (Y0 )
dY0 dx
and the equation becomes
′ 1 d 1
u0 (Y0 )u0 (Y0 ) − eY0 = −2 =⇒ (u2 ) = eY0 − 2 =⇒ u20 = eY0 − 2Y0 + B.
2 dY0 0 2
(An alternative way of getting this is to multiply the original differential equation by Y0′ and integrate). Thus we have
2
1 d
Y0 = B − 2Y0 + eY0 . (237)
2 dx
We can determine the constant B if we look at how the expansions in the previous examples matched. For example, from
Fig.16 it is seen that for the boundary-layer solution Y0′ → 0 and Y0 → y0 (0) as x → ∞. Assuming this monotonic convergence
occurs in the present example, then, since y0 (0) = ln(2), we get B = 2 [−1 + ln(2)]. Now, solving (237) for Y0′ we obtain
q
Y0′ = ± 2(B − 2Y0 + eY0 ).
To determine which sign to take, note that Y0 (0) < y0 (0). Thus, assuming the boundary-layer solution increases from
Y0 (0) = 0 to Y0 (∞) = y0 (0) = ln(2), we take the + sign in the above expression. Doing this, separating variables, and then
integrating yields the following result:
Z Y0
ds
1/2
= x. (238)
0 [2(B − 2s + es )]
This is the solution of the O(1) problem in the boundary layer, and it defines Y0 implicitly in terms of x. It is important
to note that the assumptions that were made to derive this result, such as Y0′ → 0 and Y0 → y0 (0) as x → ∞, hold for
this solution. It is also worth mentioning that another approach to matching in a case like this one involves using the phase
plane. Extensive use of this procedure can be found in O’Malley (1976,1991).
It is possible to have multiple layers in a problem and have them occur on the same side of the interval. An example of this
arises in the problem
ε3 y ′′ + x3 y ′ − εy = x3 , for 0 < x < 1, (239)

where
y(0) = 1 and y(1) = 3. (240)

The outer expansion is y ∼ x+ 2, and this holds for 0 < x ≤ 1. To investigate what is happening near x = 0, set x = x/εα ;
hence
d2 d
ε3−2α Y + ε2α x3 Y − εY = ε3α x3 . (241)
dx2 dx
There are two distinguished limits (other than the one for the outer region). One comes when 1 ∼ 3 , in which case
α = 1. This is an inner-inner layer. The other balance occurs when 2 ∼ 3 , in which case α = 21 . This is an inner layer
(see fig. 5.6). Note that the equation for the inner-inner layer does not have any terms in common with the equation for the
outer region but the inner layer equation shares terms with both. We will not investigate this any further here.
1 If y is missing from an ODE we use a different substitution to reduce the order. Consider y ′′ (x) + y ′ (x) = f (x). Let u(x) = y ′ and we reduce
to the first order ODE u′ (x) + u(x) = f (x).
56
x=0 inner outer
inner-inner
x
1/2
ε ε
Figure 21: A problem with an inner-inner layer and an inner layer.
5.6 Interior Layers

The rapid transitions in the solution that are characteristic of a boundary layer do not have to occur only at the boundary.
When this happens, the problems tend to be somewhat harder to solve simply because the location of the layer is usually not
known until the expansions have been matched. However, the expansion procedure is essentially the same as in the previous
examples. To illustrate how the method works, consider the problem
εy ′′ = yy ′ − y, for 0 < x < 1 (242)

where
y(0) = 1, y(1) = −1. (243)

The appropriate expansion in this region is the same as usual, in other words,
y(x) ∼ y0 (x) + . . . . (244)

From (242) one finds that
y0 y0′ − y0 = 0, (245)
and so either y0 = 0 or else
y0 = x + a, (246)
where a is an arbitrary constant. The fact that we have two possible solutions means that the matching might take somewhat
longer than before because we will have to determine which of these solutions matches to the inner expansion.
5.6.2 Steps 2 and 3: Interior Layer and Matching

Generally, when one first begins trying to solve a problem, it is not known where the layer(s) is. If we began this problem
as we did the previous two and assumed there is a boundary layer at either one of the endpoints, we would find that the
expansions do not match. This is a lot of effort for no results, but fortunately there is a simpler way to come to the same
conclusion. To illustrate how, suppose it is assumed that there is a boundary layer at x = 0 and (246) is the outer solution.
This situation is shown in Fig.22. If the solution behaves like the other example problems we have examined, then near x = 0
it would be expected that y ′ < 0 and y ′′ > 0. This means that εy ′′ > 0 but y(y ′ − 1) is both positive and negative in the
boundary layer. Therefore it is impossible to satisfy the differential equation (242) and have a boundary layer as indicated
in Fig.22. It is possible to rule out a boundary layer at x = 1 in the same way. It should be pointed out that these are only
plausibility arguments and they do not prove anything. What they do is guide the analysis and hopefully reduce the work
necessary to obtain the solution. However it is recommended as an exercise to see how this sort of argument might work
on a nonconvex (or nonconcave) boundary layer. (Recall a function is said to be convex when it is convex downwards i.e.,
concave upwards.)
Based on the above observations, we investigate the possibility of an interior layer. This is done by introducing the
interior-layer coordinate
x − x0
x= , (x = x0 + εγ x); Y (x) = y(x) (247)
εγ
where 0 < x0 < 1. The location of the layer, x = x0 , is not known and will be determined later. Actually, the possibilities of
either x0 = 0 or x0 = 1 could be included here, but we will not do so. Also, note that since 0 < x0 < 1 then there are two
outer regions, one for 0 < x < x0 and another for x0 < x < 1. Now, substituting (247) into (242) yields
57
y(x)
1
1 x
−1
Figure 22: Cartoon of solution with a (convex) boundary layer at x = 0 and a linear outer solution as in (246).
ε1−2γ Y ′′ = ε−γ Y Y ′ − Y. (248)

The distinguished limit here occurs when γ = 1 (i.e., the first and second terms balance). Also, as in the previous examples,
we are using Y to designate the solution in the layer. Expanding the interior-layer solution as
Y (x) ∼ Y0 (x) + . . . , (249)

it follows from (248) that Y0 satisfies the autonomous ODE:
Y0′′ = Y0 Y0′ . (250)

We can solve this by introducing the substitution:
dY0 du dY0 d2 Y0
u(Y0 ) = =⇒ = (251)
dx dY0 dx dx2
and reducing to a first order ODE.
An alternative trick is to write the equation as:
1 1 dY0 dx dY0 −dY0 dY0
Y0′′ = ( Y02 )′ ⇒ Y0′ = Y02 + α ⇒ = = = + (252)
2 2 (Y02 + β 2 ) 2 (Y0 + β)(Y0 − β) 2β(Y0 + β) 2β(Y0 − β)
where α, β are arbitrary real constants and we have assumed that α > 0. To ensure this we write α = β 2 . Integrating and
simplifying, we find that:
1 − DeBx
Y0 = B , (253)
1 + DeBx
where B and D are arbitrary constants. Other solutions of (250) are
Y0 = B tan(C + Bx/2)
obtained by assuming that α < 0 and
2
Y0 = .
C−x
assuming that α = 0.
The existence of multiple solutions is characteristic of nonlinear equations. This makes the problem interesting, but it
also means that the matching procedure is not as straightforward as it was for the linear equations studied earlier. For those
linear problems, we relied on being able to find the general solution in each region and then determining the constants by
matching. For nonlinear problems, the concept of a general solution has little meaning, and because of this it can sometimes
be difficult to obtain a solution that is general enough to be able to match to the outer expansion(s).
Of the solutions to (250), the one given in (253) is capable of matching to the outer expansions as x → ±∞. Again, it
should be remembered that the working hypothesis here is that 0 < x0 < 1 (see Fig.23). Thus, the outer expansion for
0 < x < x0 should satisfy y(0) = 1. From this it follows that
58
y0 = x + 1, for 0 ≤ x ≤ x0 . (254)
Similarly, the outer region on the other side of the layer should satisfy the boundary condition at x = 1, and this yields
y0 = x − 2, for x0 < x ≤ 1. (255)

Now, for (253) to be able to match either (254) or (255) it is going to be necessary that both B and D in (253) are
nonzero (in fact, without loss of generality, we will take B to be positive). The requirements imposed in the matching are
very similar to those obtained for boundary layers. (Exercise: Do the matching using van Dyke). In particular, we must
have that Y0 (∞) = y0 (x+ −
0 ) and Y0 (−∞) = y0 (x0 ).
From (253) and (254) we get
B = x0 + 1,
and from (253) and (255) we have
−B = x0 − 2.
3
Solving these equations, one finds that B = 2 and x0 = 21 .
Note It remains to determine D in (253), and the only requirement on this constant, so far, is that it is nonzero. It might
seem strange that we have seemingly satisfied everything for the first-order problem, including matching, yet we still have
not completely determined the first term in the expansion. What we are missing is the global structure that is imposed by
the problem. This can be missed when patching the solution together like we have done using inner and outer layers and
simple power series expansions. One way to resolve this situation can be found by looking a little closer at Fig.23. It suggests
the solution has a symmetry with respect to the point where the curve crosses the x-axis, which appears to be at x = 1/2. To
prove this, suppose the boundary conditions are more general than those given in (243). In particular, suppose y(0) = a and
y(1) = b. The solution can then be written as y = f (x, a, b). Now, if we change variables and let s = 1 − x (which produces
a flip around x = 1/2) and z = −y (which produces a flip around y = 0), then
εz ′′ = zz ′ − z, for 0 < s < 1,

where x(0) = −b and x(1) = −a. This is the same problem as in (242), and therefore the solution is z = f (s, −b, −a).
Transforming back yields y = −f (1 − x, −b, −a). The solution is unique, and so f (x, a, b) = −f (1 − x, −b, −a). This odd-
looking relationship is the symmetry condition we are looking for, and we now take x = 1/2, a = 1, and b = −1. In this case,
we find that f ( 21 , 1, −1) = 0, and therefore y( 21 ) = 0. With this it follows that D = 1 in (253), and

3 3x
Y (x) ∼ − tanh . (256)
2 4
The symmetry argument used to determine the unknown constant is relatively easy to carry out once the symmetry is
known. However, finding a symmetry can be difficult, if not impossible. Fortunately, there are other ways to resolve this
problem e.g. Grasman and Matkowsky (1977). Other methods include a phase-plane analysis (O’Malley 1976, 1991), an
extension of the variational approach that eliminates exponentially small eigenvalues of an associated linearized problem
(Keth et el., 1987; Ward, 1992), and the use of exponentially small terms in the asymptotic expansions (Lange, 1983).
Because all of these have limitations in terms of their applicability, there is still work to be done in resolving this problem.

It can be more difficult to construct a composite expansion when the outer solutions are discontinuous across the interior
layer like they are in this problem. What is done is to find one for the interval 0 ≤ x ≤ x0 and then another for x0 ≤ x ≤ 1.
As it turns out, for this example the expansions on either side are the same, and the result is
3
y ∼x+1− , for 0 ≤ x ≤ 1. (257)
1+ e−3(2x−1)/4ε
This composite expansion is graphed in Fig.24 to illustrate the nature of the interior layer. The rapid transition from
one outer solution to the other is typical of what is sometimes called a shock solution Also shown in Fig.24 is the numerical
solution, and it is clear that the composite expansion is in close agreement with the numerical solution.
The interior layer of the previous example consists of a monotone transition from one outer solution to another. Other
types of interior layers are also seen, including spikes and nonmonotone transition layers. Examples of these are shown in
Fig.25. There are also some very unusual interior layers. For example, O’Malley (1976) gives an example where there is a
spike at each rational point on the interval −1 < x < 1. An analysis of nonmonotone interior-layer problems can be found in
DeSanti (1987).
59
y(x)
0
1 x
-1
Figure 23: The linear functions which make up the outer expansion and the interior layer solution connecting them.
2
Numerical
Solution
1 Composite
Solution
Solution
-1
-2
0 0.2 0.4 0.6 0.8 1.0
x - axis
Figure 24: Composite expansion (257) and the numerical solution.
1.2
0.9
0.6
0.3
0
a
1.0
0.6
0.2
-0.2
-0.6
b 0 0.2 0.4 0.6 0.8 1.0
x - axis
Figure 25: Examples of nonmonotone interior layers. In the first there is a spike generated by the function y(x) =sech2 [(x −
1 1 1 2
2 )/ε], ε = 0.01. In the second there is a nonmonotone transition layer generated by y(x) = 2 tanh[(x − 2 )/ε]+sech [(x −
1
2 )/ε].
60
5.7 Corner Layers
One of the distinguishing features of the problems we have studied in this chapter is the rapid changes in the solution in the
layer regions. The problems we will now investigate are slightly different because the rapid changes will be in the slope, or
derivatives, of the solution, and not in the value of the solution itself. To illustrate this, we consider the following problem:

′′ 1
εy + x − y ′ − y = 0, for 0 < x < 1, (258)
2
where
y(0) = 2 (259)
and
y(1) = 3. (260)
′
Incidentally, the coefficient of the y term in this equation changes sign in the interval. For this reason, the equation is said
to have a turning point at x = 1/2. By itself this means little, but as we will see it does play a role in this example by
preventing boundary layers at the endpoints of the interval. This also happened with the nonlinear problem in the previous
section, which ended up having an interior layer.

The solution in this region is expanded in the usual power series as follows:
y(x) ∼ y0 (x) + εy1 (x) + . . . (261)

From (258) one then finds that

dy0 y0 1 1
= =⇒ ln y0 = ln(x − ) + b =⇒ y0 = a x − , (262)
dx x − 21 2 2
where a is an arbitrary constant. As usual, we must satisfy two boundary conditions with only one integration constant.
5.7.2 Step 2: Corner Layer

We begin by determining whether there are boundary layers. These can be ruled out fairly quickly using the plausibility
argument presented in the previous section. For example, if there is a boundary layer at x = 0 and (262) is the outer solution,
then we have a situation similar to that shown in Fig.22. In the boundary layer y ′′ > 0, y ′ < 0 and there is a portion of the
′′ 1 ′
curve where y < 0. This means that εy > 0 > − x − 2 y + y. It is therefore impossible to satisfy the differential equation
(258) and have a boundary layer as indicated in Fig.22. Using a similar argument, and the fact that the coefficient of the y ′
term changes in the interval, one can also argue that there is not a boundary layer at the other end.
It therefore appears that there is an interior layer, and to investigate this we introduce the stretched variable
x − x0
x= , (x = x0 + εα x). (263)
εα
where we expect that x < 0(x < x0 ) and x > 0(x > x0 ).
With an interior layer there is one outer solution for 0 ≤ x < x0 and another one for x0 < x ≤ 1. Using boundary
conditions (259) and (260) and the general solution given in (262), we have that the outer solution is

−4 x − 12 for 0 ≤ x < x0 (left)
y∼ (264)
6 x − 12 for x0 < x ≤ 1 (right)
Taking x0 = 12 , we find that the outer solution is continuous, giving us the situation shown in Fig.26. The outer solution is
continuous but is not differentiable at x = 21 . It is for this reason that we have what is called a corner region, or derivative
layer, at x0 = 21 .
Now, substituting (263) into (258), and letting Y designate the solution in this region, we get
ε1−2α Y ′′ + xY ′ − Y = 0. (265)
1
The distinguished limit here occurs for α = 2. Unlike what we have assumed in the previous examples, we now take
Y ∼ Y−1 (x) + εγ Y0 + . . . . (266)

In fact we will take Y−1 = 0 for this particular problem. A better way of doing this would be to perform a rescaling. (This
will be borne out by the matching and indicates that in the boundary layer y = O(εγ )). Also, the multiplicative factor εγ is
61
y(x)
3
2
1/2 1 x
Figure 26: Cartoon of the corner layer of the outer solution (264). The solution from the corner layer region should provide
a smooth transition between these linear functions.
needed to be able to match to the outer solution (the constant γ will be determined from the matching). So, substituting
this into (265) yields
Y0′′ + xY0′ − Y0 = 0, for − ∞ < x < ∞. (267)

We are in the unfortunate situation of having to solve the full problem to be able to determine what is going on in this layer.
If we can solve this problem, one might wonder why it is necessary to bother with an asymptotic approximation in the first
place. There is no arguing with this sentiment, and this is one of the situations where perturbation methods show some of
their limitations. However, there are several comments that should be made. First, it is sometimes easier to solve a problem
on an infinite, or semi-infinite, interval as occurs in an interior or boundary-layer problem. Second, even if we cannot solve
the layer problem in closed form, it is usually possible to say a great deal about the structure of the solution.
As it turns out, it is possible to solve (267) in closed form by using power series methods. However, a simpler way is to
notice that Y0 = x is a solution, and so, using the method of reduction of order, one finds that the general solution is
" Z #
x
2 2
Y0 = Ax + B −e−x /2
−x e−s /2
ds . (268)
0
5.7.3 Step 3: van Dyke matching

We have two outer solutions (264) to O(1) ≡ ∆ and an inner solution εγ Y0 where Y0 is given by (268). Thus the inner
solution is known to O(εγ ) ≡ δ . The idea here is to match the first (left) outer solution to the corner layer solution and
then to match the corner layer solution with the (right) outer solution (see fig.26).
1
We write the left hand outer solution in terms of the inner variable (x = 1/2 + ε 2 x) and we expand this for x = O(1) for
small ε
1
(left) (y0 )i = −4(ε 2 x) (269)
where we expand this as far as O(εγ ) ≡ δ and it is sensible to take γ = 1/2 in order to match
√ the solutions.
Similarly we write the corner solution in the form (268) (N.B. the corner solution is εY0 not just Y0 ) in terms of the
outer variable (x = ( x−1/2
√
ε
)) and we expand this for x = O(1) and x < 1/2 for small ε
" Z ( x−1/2 # r
)
√ √ (2x − 1)2
√
o
ε
−s2 /2 π
( εY0 ) = A(x − 1/2) + B − ε exp(− ) − (x − 1/2) e ds ∼ A(x − 1/2) + B(x − 1/2) (270)
8ε 0 2
Z −∞
2 pπ
where we expand as far as ∆ = O(1). Note that we have used the fact e−s /2
ds = − 2. (Watch the uppermost limit
0
in the integral). We now wish to compare (269) and (270). We first write the former in terms of x rather than x and we see
that we require
r
π
−4(x − 1/2) = A(x − 1/2) + B(x − 1/2) (271)
2
62
so we require
r
π
−4 = A + B . (272)
2
1
Matching the corner solution with the right outer solution we write the latter in terms of the inner variable (x = 1/2 + ε 2 x)
and we expand this for x = O(1) for small ε to get
1
(right) (y0 )i = 6(ε 2 x) (273)
√ x−1/2
while we write the corner solution ( εY0 ) in terms of the outer variable (x = ( √ε )) and we expand this for x = O(1) and
x > x0 (= 21 ) for small ε
" Z ( x−1/2 # r
)
√ √ (2x − 1)2
√
o
ε
−s2 /2 π
( εY0 ) = A(x − 1/2) + B − ε exp(− ) − (x − 1/2) e ds ∼ A(x − 1/2) − B(x − 1/2) (274)
8ε 0 2
Z ∞
2 p
where we have used e−s /2 ds = π2 . Writing (273) in terms of x and comparing the last two equations we see that
0
r
π
6(x − 1/2) = A(x − 1/2) − B(x − 1/2) (275)
2
so that
r
π
6=A−B (276)
2
p
From (272) and (276) w see that A = 1 and B = −5 2/π.

Using van Dyke matching we first construct a composite solution in [0, 1/2] and we find by adding (264) and (270), and
subtracting (269) (all written in terms of x) we get
" Z ( x−1/2 #
p )
√ −(2x − 1)2
√
ε
−s2 /2
(x − 1/2) − 5 2/π − ε exp( ) − (x − 1/2) e ds (277)
8ε 0
while on [1/2, 1] we in fact obtain exactly the same expression for the composite solution.
This function is shown in Fig.27 and is clearly in very good agreement with the exact solution.
Note on reduction of order Suppose we have a linear ODE of order n of the form L[y] = 0 and y1 is a known solution.
Then we can seek a second solution by putting y2 = u(x)y1 . This yields an equation M [u] = 0 which is linear homogeneous
and of order (n − 1) in terms of the dependent variable v(x) ≡ du dx . The idea is that it should be easier to find v(x̄), hence
u(x̄), and finally a second independent solution of the original ODE u(x̄)y1 (x̄) ≡ y2 (x̄).
In the above example we let Y0 = u(x)x so that
Y0′ = xu′ + u; Y0′′ = xu′′ + 2u′ . (278)

Now (267) yields
xu′′ + u′ (2 + x2 ) = 0 ⇒ xv ′ + v(2 + x2 ) = 0; (v ≡ u′ ) (279)

This is a first order linear equation for v(x) which can be solved (in this instance) by separating the variables. One finds
that
du const
v= = 2 2 (280)
dx x exp( x2 )
and this may be integrated using integration by parts to get
x2 Z ! Z
e− 2 −1 −x2 2
−1 − x2 −x2
u=A − x d(e 2 ) =A x e + e 2 dx (281)
x
and so a second independent solution of the ODE is

Z
2 −x2
− x2
y2 (x̄) ≡ xu(x̄) = A e +x e 2 dx (282)
which leads to (268).
63
Intermediate variable matching and composite solution In the examples from the previous two sections, the layer
analysis and matching were carried out in a single step. This is not done here because the matching in this problem is slightly
different and is worth considering in more detail. To do the matching, we introduce the intermediate variable
x − (1/2)
xη = , (283)
εk
where 0 < k < 21 . Rewriting the outer solution (264) in this variable yields

−4εk xη for xη < 0
y∼ . (284)
6εk xη for 0 < xη
Also, using the fact that
Z ∞ r
−s2 /2 π
e ds = ,
0 2
it follows from (268) that
p
εγ+k−1/2 xη (A + B p π2 for xη < 0
Y ∼ (285)
εγ+k−1/2 xη (A − B π2 for 0 < xη
To be able to match (284) and (285), we must have γ = 21 . In this case,
r
π
A+B = −4
2
and
r
π
A−B = 6,
2
p
from which it follows that A = 1 and B = −5 2/π.
Even though the situation is slightly more complicated than before, the construction of a composite expansion follows
the same rules as in the earlier examples. For example, for 0 ≤ x ≤ 21 ,
( " Z x #)
1 1/2
p −x2 /2 −s2 /2
y ∼ −4 x − + ε x 1 − 5 2/π −e − e ds + 4ε1/3 xη (286)
2 0
( r " Z (2x−1)/(2√ε) #)
1 2 −(2x−1)2 /(8ε) −s2 /2
= x− 1+5 e + e ds .
2 π 0
1
One finds that this is also the composite expansion for 2 ≤ x ≤ 1. Thus, (286) is a composite expansion over the entire
interval.
5.8 Partial Differential Equations

The subject of boundary and interior layers and how they appear in the solutions of partial differential equations is enormous.
We will examine two examples that lend themselves to matched asymptotic expansions.
5.8.1 Example
A Parabolic PDE Problem To illustrate the application of boundary-layer methods to parabolic PDEs, we consider the
problem of solving
ut + uux = εuxx, for − ∞ < x < ∞ and 0 < t, (287)

where u(x, 0) = f (x). It is assumed here that f (x) is smooth and bounded except for a jump discontinuity at x = 0.
Moreover, f ′ (x) ≥ 0 for x 6= 0 and f (0− ) > f (0+ ).
It turns out that if f (x) is ever a decreasing function of x, this leads to wave steepening and wave-breaking (if we allow
this). For this reason in much of what follows f (x) is a simple step function 1 − H(x), where H(x) is the Heaviside function.
In such cases f (x), the initial profile, is a wave that is about to “break”.
The nonlinear advection-diffusion equation in (287) is known as Burgers’ equation. It has become the prototype problem
for studying shock waves and for the use of what are called viscosity methods for finding smooth solutions to such problems.
In what follows, we will concentrate on the constructive aspects of finding asymptotic approximations to the solution. The
theoretical underpinnings of the procedure can be found in Il’in (1992) and its extension to systems in Goodman and Xin
(1992).
64
3 Exact Solution
Composite Expansion
2 Outer Expansion
Solution
1
-1
0 0.2 0.4 0.6 0.8 1.0
x - axis
Figure 27: Composite expansion (286), the exact solution, and the outer expansion (264). (ε = 0.01).
Outer Expansion The first step is to find the outer solution, and to do this we assume u ∼ u0 (x, t) + . . . . Substituting
this into (287) produces the first-order hyperbolic equation
∂t u0 + u0 ∂x u0 = 0, u0 (x, 0) = f (x). (288)

We solve this PDE problem using the method of characteristics. Note that the full PDE in (287) is parabolic (it is an
advection diffusion equation in fact). When we omit the O(ε) terms, the reduced PDE becomes hyperbolic. Because a
parabolic equation does not have such real characteristics, we sometimes use the terminology sub-characteristics in problems
of this type (where the reduced problem is hyperbolic but the full problem is not). This is analogous to a singular perturbation
ODE problem changing order. We first parameterise the initial condition to the form:
x = α, t = 0, u0 (α, 0) = f (α), −∞ < α < ∞. (289)

The characteristic curves are defined by:
dx
= u0 (290)
dt
and along these curves, the PDE reduces to the ODE:
∂u0 ∂u0 dx du0 (x(t), t)

+ =0⇒ = 0. (291)
∂t ∂x dt dt
This has solution u0 = constant (along a characteristic) and from the initial conditions we see that:
u0 = f (α) (292)
so the characteristics are given by:
dx = f (α)dt ⇒ x = f (α)t + A (293)

Using the boundary conditions we find the integration constant in terms of α and we see that A = α so the characteristics
are x = f (α) + α and the solution in parametric form is:
x = f (α)t + α, u0 = f (α) (294)

From the first equation we see that α = x − f (α)t = x − u0 t so we can eliminate the parameter α to find the solution (in
implicit form):
u0 = f (x − u0 t). (295)
65
Note In finding the solution in parametric form, we used the fact that u0 is constant along a characteristic curve. Once
the solution has been established, we see that it defines u0 in some region of x − t space and we can interpret the solution
more generally via:
u0 = f (α) with α = α(x, t) defined implicitly by α = x − f (α)t.
Here we are using the fact that provided the characteristics have not crossed, it is clear from the characteristics diagram that
any arbitrary point (x, t) corresponds to a unique value of α.
Referring to fig. 28, because the local wave speed is u0 , if we start with an initial condition f (x) as shown, the wave will
tend to steepen at the front (and eventually break) at which point the solution in multi-valued. Usually this is physically
unacceptable.
More formally, in moving from the parametric solution to the explicit solution, we need to invert the transformation
x(α), t(α), and to express it in the form α(x, t) i.e., we have to solve for α as a function of (x, t). In the present instance the
inversion is easy but this is not always the case. There are also situations where a double parameterisation is required and
these can be even more complicated. (See, for example, page 29 An Introduction to PDES, Pinchover and Rubinstein).
u
u
x
u is
multi-valued
Figure 28: Wave steepening and breaking for the “inviscid” Burgers’ equation i.e., (287) with ε = 0. If the initial condition
f (x) is decreasing on any interval, this will always lead to wave steepening.
We can analyze this phenomenon more closely by examining the characteristic curves. We can graph these in xt space.
The characteristics for this equation are the straight lines x = α + f (α)t, and the solution u0 (x, t) is constant along each
of these lines (see Fig.29). Therefore, given a point (x, t), then u0 (x, t) = f (α), where α is determined from the equation
x = α + f (α)t. This construction succeeds if the characteristics do not intersect and they cover the upper half-plane. For
us the problem is with intersections which correspond to the “wave” breaking and (inadmissable) multivalued solutions.
Consequently, based on the nonlinear hyperbolic nature of the reduced problem in (288), we will assume there is a single
(smooth) curve x = s(t) across which the outer solution is discontinuous (see fig 29). The solution will jump in value as we
move across this curve. As this curve is in xt space, it represents the location of a moving discontinuity. The precise location
of this curve will be determined from the analysis of the transition layer that connects these two outer regions.
To summarise we assume that our outer solution has developed to the point where it is about to form a discontinuity
(shock) moving at speed s(t), to be determined, with value u− +
0 on the left and u0 on the right. We wish to examine the
details of the solution in the vicinity of the shock (the shock structure) by rescaling the governing equation (287).
5.8.2 A special case

In order to simplify the discussion we will now specialise to the case where the shock (or discontinuity in the solution) exists
from t = 0. Specifically, consider the problem (287) with initial conditions given by the step function:
u(x, 0) = u− +
0 , x < 0; u(x, 0) = u0 , x ≥ 0 (296)
so that the discontinuity is already present in the initial condition (we assume >u−
0 u+
0 ).
The discontinuity assumed to be
moving at speed ds(t)/dt is to be determined in the following inner layer analysis. (Note that the following analysis would
66
t
x=s(t)
Figure 29: The straight lines are the characteristics for the reduced problem (288). Because f has a discontinuity at x = 0,
′
and f ≥ 0 for x 6= 0, these lines intersect along the shock curve x = s(t).
still hold in the general case). The solution at this point is thus a step function with the step located at s(t), with s(t) to
be determined. In this case, the characteristic curves are give by figure 33 and figure 34 shows the insertion of the shock
location. The slope of the characteristics on either side of x = 0 is given by (290). Thus characteristics emanating from
x < 0 have slope 1/u− 0 , while those emanating from x > 0 have slope 1/u0 .
+
Inner Expansion In order to rescale (i.e., rebalance the terms in the governing equation) in the vicinity of the shock, we
define a stretched inner layer coordinate which moves with the speed of the moving shock i.e., s′ (t).Thus it seems reasonable to
put x = s(t) + εα x. Let U (x, τ ) denote the solution of the problem in this layer. We wish to transform from u(x, t) → U (x, τ )
where x = [x − s(t)] /εα , τ ≡ t. Using the chain rule we see that
∂u ∂U ∂x ∂U ∂τ −s′ (t) ∂U ∂U
≡ + = + (297)
∂t ∂x ∂t ∂τ ∂t εα ∂x ∂τ
and
∂u ∂U ∂x ∂U ∂τ 1 ∂U
≡ + = α +0 (298)
∂x ∂x ∂x ∂x ∂x ε ∂x
so we see in operator notation that
∂ −s′ (t) ∂ ∂ ∂ 1 ∂ ∂2 1 ∂2
≡ + ; ≡ ; ≡ (299)
∂t εα ∂x ∂τ ∂x εα ∂x ∂x2 ε2α ∂x2
So (287) takes the form
∂t U − ε−α s′ (t)∂x U + ε−α U ∂x U = ε1−2α ∂x2 U (300)

where we have chosen to rewrite τ as t for the time variable.
Balancing terms in this equation, one finds α = 1. Thus, the appropriate expansion of the solution is U ∼ U0 (x, t) + . . . ,
and substituting this into (300) yields at leading order
−s′ (t)∂x U0 + U0 ∂x U0 = ∂x2 U0 .

It is important to note that the explicit time derivative ∂t U in (300) disappears at leading order and the equation is a
pseudo PDE. This suggests that in the shock at leading order the profile is quasi-steady i.e., that the characteristic time
scale for the shock is much shorter than that for the outer solution: the shock reacts very quickly (instantaneously at leading
order) to any changes which occur in the outer region. These (time dependent) effects will be incorporated via matching
with the outer solution. (In fact it is easy to see that the characteristic O(1) time scale T for the shock is defined by t = εT
and if we had used x, T as our inner variables, we would have included an explicit time derivative term ∂T U in the inner
equation).
One way to integrate this (autonomous in x) pseudo PDE is to let ∂x U0 = v(U0 ) so that ∂x v(U0 ) = dv(U
dU0
0 ) ∂U0 dv(U0 )
∂x = U0 dU0 .
′
Another possibility is to use the substitution g(x, t) =s (t) − U0 so that the equation becomes
67
1
g∂x g + ∂x2 g = 0 ⇒ ( g 2 )x + ∂x2 g = 0 (301)
2
which has a first integral
1 2
g + ∂x g = A(t) (302)
2
where A(t) is an integration function. The boundary conditions to be used come from matching with the outer solution.
They are
lim U0 = u−
0 and lim U0 = u+
0, (303)
x→−∞ x→∞
In terms of g(x, t) these conditions become
lim g = s′ (t) − u−
0 and lim g = s′ (t) − u+
0 (304)
x→−∞ x→∞
2 2
which also means that ∂x g → 0 as x → ±∞. From (302) it follows that A(t) = 1
2 s′ (t) − u−
0 = 1
2 s′ (t) − u+
0 and so
+2 −2
1
2 (u0 − u0 ) [ 12 u20 ]+
− 1 + 1 − − 2
s′ (t) = = + = u0 + u−
0 ; A(t) = (u0 u0 ) . (305)
u+ −
0 − u0 [u0 ]− 2 8
where the square bracketed terms denote a jump in a quantity across the shock.
The differential equation for s(t) in (305) determines the position of the shock and is known as the Rankine-Hugoniot
condition. Its solution requires an initial condition, and because of the assumed location of the discontinuity in f (x) we take
s(0) = 0.
To complete the analysis of the first-order problem in the shock layer, we separate variables in (302) writing for clarity
1 − − 2
A(t) = (u u ) = β 2 /2 > 0.
8 0 0
Thus
1 2 2dg
g + ∂x g = β 2 /2 =⇒ 2 = dx (306)
2 (β − g 2 )
R dg
But (β 2 −g2 ) = 1
β tanh−1 g
β so
β
g = β tanh (x + k) (307)
2
and in terms of U0 this gives:
β
U0 (x, t) = s′ (t) − g = s′ (t) − β tanh
(x + k) (308)
2
Note that the tanh function (see fig. 5.8.2) has the correct qualitative behaviour to resolve the discontinuity in the shock.
Note that k = k(t) is another integration function which arises from integrating the pseudo PDE. In principal it can be
obtained from matching with the outer solution but we note that it merely corresponds to a translation of the shock in the
x direction and can be taken to be zero here because s(0) = 0.
Comment Here, we will not be further concerned with the details of the matching. We will content ourselves with the
fact that the inner solution shows the correct qualitative behaviour (via the tanh function) allowing the outer solution to
step down (smoothly) from u− +
0 to u0 .
To demonstrate the accuracy of the asymptotic approximation, we show in Fig.32 the approximation along with the exact
solution at two different values of t and ε. For ε = 0.1 the two solutions are so close that they are essentially indistinguishable
in the graph. They are even in reasonable agreement when ε = 1, although the differences are more apparent for the smaller
value of t. The sharpening of the shock layer as ε decreases is also evident in these figures. This is expected because we
found that the width of this layer is O(ε).
68
[ptb]
1
0.5
–3 –2 –1 0 1 2 3
–0.5
–1
Figure 30: tanh(x) and − tanh(x).
t x=s(t)
x- x+ x
Figure 31: Schematic showing the shock and two characteristics that intersect on the shock.
69
1.0
0.8
u(x,t)
0.6
t = 10 t = 20
0.4
0.2
0
-5 0 5 10 15 20
x - axis
1.0
0.8
t = 20
u(x,t)
0.6
0.4
t = 10
0.2
0
2 4 6 8 10 12
x - axis
Figure 32: The exact solution of Burger’s equation (dashed curves) and the asymptotic approximation (solid curves) where
φ(x) = 1 if x < 0 and φ(x) = 0 if x > 0. The solutions are compared at t = 10, t = 20 with ε = 1 and ε = 0.1 respectively.
Comment Consider the outer problem again
∂t u0 + u0 ∂x u0 = 0 (309)
with boundary conditions (296) in the specific case where u− +
0 = 1, u0 = 0 or f (x) = 1, x < 0, f (x) = 0, x ≥ 0. Thus from
(293) the characteristics are vertical if x > 0 and they are straight lines with slope 1 if x < 0 (see fig. 33). Immediately at
t = 0 one sees that the characteristics collide and a contradiction is implied since u0 must be a constant along a characteristic.
(The solutions might be considered to be multi-valued, but this is usually not physically acceptable). One way to avoid this
impasse is to insert a curve (in fact for this simple problem this curve turns out to be a straight line x = s(t) = mt (m > 0))
along which the initial discontinuity at x = 0 is carried. This changes the characteristic diagram to fig. 34. For x > mt we
can take u = 0, and for x < mt we can take u = 1, thus giving a solution to the PDE on both sides of the discontinuity.
The only question is the choice of m = ds/dt (which in this case is the shock speed). It appears that any curve starting
at the origin of xt space and continuing through the first quadrant will solve the issue of multivaluedness. For the present
problem this means that any m ≥ 0 will give rise to such a curve in the first quadrant and this also satisfies the initial
condition. Should we give up uniqueness for this problem? Is there some other solution that we have not discovered? Or is
there some special choice of m? As we saw above, the inner problem gives us a way of choosing the shock speed ds/dt and
hence choosing the slope m i.e., ds/dt is defined by(305).
Comment We make one further comment on the jump conditions across the shock and the shock speed. Consider again
the outer equation (288) or (309). This can also be written in the form
1 2
u
∂t u0 + ∂x (q) = 0; q ≡
2 0
This equation is said to be in conservation form: evolution equations like this are generally derived from a conservation law
(the terminology “weak” formulation is also used), here of the form
Z B
d 1 1
u0 dx = q(A, t) − q(B, t) = u0 (A, t)2 − u0 (B, t)2 , (310)
dt A 2 2
Physically this means that increase in the quantity of u0 between A and B is determined by the flow in at A and out at B.
The deduction of the point form (288) from (310) required the additional assumption that u was continuously differentiable;
however, it is possible to satisfy (310) at a point of discontinuity of u. Suppose u is discontinuous at the smooth curve
x = s(t), that u is continuously differentiable for u > s(t) and u < s(t) and that u and its derivatives have finite one-sided
limits as x → s(t)− and x → s(t)+ . Choose A < s(t), B > s(t). Denote the jump in a quantity q across the shock by
1 2
[q]+ + −
− ≡ q(x = s(t) , t) − q(x = s(t) , t), q = u .
2 0
Write (310) in the form:
70
t
u=1 u=0 x
Figure 33: Characteristics for the initial value problem with u− +

0 = 1, u0 = 0.
x=mt
u=1
u=0
x
u=1 u=0
Figure 34: Insertion of the line x = mt along which the discontinuity is carried. Note that any curve in the second quadrant
passing through the origin can be used to split the characteristics but only one of these curves is the correct one.
71
Z s(t) Z B
d d 1 1
u0 dx + u0 dx = u0 (A, t)2 − u0 (B, t)2 .
dt A dt s(t) 2 2
Recall Leibniz’ rule Z Z
g(t) g(t)
d ∂u dg(t) df (t)
u(x, t)dx = dx + u(g(t), t) − u(f (t), t)
dt f (t) f (t) dt dt dt
for differentiating under the integral sign. Applying this we obtain:
Z s− (t) Z B
∂u0 ∂u0 ds ds 1 1
dx + dx + u0 (x = s− , t) − u0 (x = s+ , t) = u0 (A, t)2 − u0 (B, t)2 ,
A ∂t s+ (t) ∂t dt dt 2 2
where u0 (s− , t), u0 (s+ , t) are the one-sided limits of u0 (x, t) as x → s(t)− and x → s(t)+ respectively. Here u0 is continuously
differentiable to the left and right of the shock. Now take the limit as B → s(t)+ , A → s(t)− : the first two terms go to zero
because the integrand is bounded and the integral of integration goes to zero. Thus we find that (310) implies the jump
condition:
ds(t) [ 1 u 2 ]+ 1 + [q]+
= 2 0+− = u0 + u−
0 =
−
dt [u0 ]− 2 [u0 ]+
−
as we already obtained in (305). It is gratifying that the analysis of the inner problem gives a shock speed condition which
is consistent with that found by analysing the outer problem in conservation form.
Comment Note finally that typically (310) will be derived from the physics of the problem under consideration (and
represents conservation of some quantity). It is a more fundamental representation of the problem than (309) in that (310)
allows solutions with discontinuities or shocks while the DE form (309) can only admit smooth solutions (sufficient smoothness
requirements are that u, q are continuously differentiable). If the solutions of a particular problem are smooth (or classical
solutions) then both forms are applicable and the DE form is usually more convenient. If shocks (i.e., discontinuities in u, q),
occur, the integrated form (310) must be used, at least in the vicinity of the shock. The integrated form of the conservation
law is sometimes referred to as a weak formulation of the problem and solutions with shocks are then termed weak or
non-classical solutions.
We thus apparently have two ways of dealing with problems of this type to resolve the multivaluedness in the solution
which arises when the characteristics cross:
• perform a rescaling of the problem in the moving shock layer of thickness ε (this gives a smoothed shock profile and
the shock speed) or
• use the weak formulation starting with the equation in conservation form (310) (this results in a solution with a
discontinuity (shock) and also gives the location and speed of this shock).
Note that for the more general conservation law:
ut + qx = 0,
an analysis similar to the above yields the shock speed condition:
ds(t) [q]+
= − .
dt [u]+
−
Of course, in the present case, q = 21 u2 .
Summary We summarise by saying that in solving Burgers’ equation with ε = 0 (sometimes called the “inviscid” Burgers’
equation), where unacceptable multivalued solutions develop, these can be resolved either by including the O(ǫ) smoothing
terms which leads to a smooth solution, or by inserting a vertical discontinuity or shock on mass conservation grounds.
Typically the results of the two approaches are summarised in fig. 35. The two results coincide as ε → 0.
In this section, we have considered elliptic and parabolic problems. Matched asymptotic expansions can also be applied to
hyperbolic problems, and several examples are worked out in Kevorkian and Cole (1981). It should also be pointed out that
the methods developed in (6) and (8) are applied extensively to hyperbolic equations.
72
Inviscid Viscous
Discontinuity Continuous
in u u
Figure 35: Two approaches for “fixing” the multivalued solutions of Burgers’ equation.
5.9 Extra problems

5.9.1 Example
The equation studied in this section is related to one that has received a great deal of attention because the solution has
been found to have some rather interesting properties. To explain the situation, consider

1
εy ′′ − x − y ′ + ky = 0, for 0 < x < 1, (311)
2
where k is a constant. This equation has a simple turning point at x = 21 . Assuming that y(0) = y(1) = 1, then the solution
is

M − k2 , 21 , 8ε
1
(2x − 1)2
y(x) = , (312)
M − k2 , 12 , 8ε
1
where M is confluent hypergeometric function (see Appendix*). It is assumed here that the value of ε is such that the
denominator in (312) is nonzero. Using the known asymptotic properties of M , one finds that for small ε and x 6= 1 − 2,
−x(1−x)/2ε
e for k 6= 0, 2, 4 . . .
y(x) ∼ .
(2x − 1)k for k = 0, 2, 4 . . .
This shows that there is a boundary layer at each endpoint for all but a discrete set of values for the constant k. When
there are boundary layers, the solution in the outer region is transcendentally small and goes to zero asε ↓ 0. What is
significant is that for k = 0, 2, 4 . . . this does not happen. This behaviour at a discrete set of points for the parameter k
is reminiscent of resonance and has become known as the Ackerberg-O’Malley resonance. Those who may be interested in
pursuing this topic are referred to the original paper by Ackerberg and O’Malley (1968) and a later study by De Groen
(1980).
5.9.2 Example: An elliptic PDE

The first example concerns finding the function u(x, y) that is the solution of the following boundary-value problem:
ε∇2 u + α∂x u + β∂y u + u = f (x, y), for (x, y) in Ω, (313)

where
u = g(x, y), for (x, y) on ∂Ω. (314)
73
Figure 36: The solution of (315) and associated contours (a = 1, b = 0, ε = 0.05). For this problem Ω is x2 + y 2 ≤ 1. The
solution outside this region can be ignored. Clearly there is a boundary layer.
The domain Ω is assumed to be bounded and simply connected, and to have a smooth boundary. The coefficients α and
β are constants with at least one of them nonzero. The functions f (x, y) and g(x, y) are assumed to be continuous.
To illustrate what a solution of (313) looks like, consider the special case where f (x, y) = a and g(x, y) = b, where a and
b are constants.Taking Ω to be the unit disk and using polar coordinates x = ρ cos ϕ, y = ρ sin ϕ, then the solution of the
problem is
X∞
1
u(ρ, ϕ) = a + In (χρ) [an sin(nϕ) + bn cos(nϕ)] /In (χ), (315)
F (ρ, ϕ) n=−∞
ρ
√
where F (ρ, ϕ) = exp 2ε2 (α cos ϕ + β sin ϕ) , χ = 2ε12 2 − 4ε2 , and In is a modified Bessel function,
Z 2π
b−a
an = F (1, ϕ) sin(nϕ)dϕ
2π 0
and
Z 2π
b−a
bn = F (1, ϕ) cos(nϕ)dϕ.
2π 0
This solution is shown in figure 36. A few observations can be made from this figure that will make the derivation of the
asymptotic approximation
√ easier to follow. First, a boundary layer is clearly evident in the solution. For example, if one
starts at x = y = −1/ 2 and then moves into the domain, the solution undergoes a rapid transition to what appears to be
an outer solution. It is also seen that the layer is not present around the entire boundary ∂Ω. Moreover, if one follows the
solution around the edge of the boundary, it’s not easy to identify exactly where the solution switches between the boundary
layer and the outer region. However, wherever it is, the transition is relatively smooth. This latter observation will be useful
later when deciding on the importance of what are called tangency points.
Before jumping into the derivation of the asymptotic expansion, it is worth noting that when ε = 0 the differential
equation in (313) reduces to a first-order hyperbolic equation. This change in type, from an elliptic to a hyperbolic equation,
has important consequences for the analysis and helps explain some of the steps that are taken next.
Outer Expansion The procedure that we will use to find an asymptotic approximation of the solution is very similar to
that used for ordinary differential equations. The first step is to find the outer solution and to do this assume
u ∼ u0 (x, y) + εu1 (x, y) + . . . . (316)

Substituting this into (313) yields the following O(1) problem:
α∂x u0 + β∂y u0 = f (x, y), in Ω, (317)

where
74
y
P
*
TL
Ω TR
x
P
Figure 37: Cartoon of the characteristic curves obtained for the outer solution. These are determined by (319) and are
directed staright lines with an orientation determined by the direction of increasing s.
u0 = g(x, y), on ∂Ωo . (318)

When constructing an asymptotic approximation, it is not always immediately clear which boundary condition, if any, the
outer solution should satisfy. In (318) the portion of the boundary where u0 satisfies the original boundary condition has
been identified as ∂Ωo . This is presently unknown and, in fact, it could turn out to be empty (Exercise). In any case, to
solve (317) we change coordinates from x, y to r, s, where s is the characteristic direction for (317). Specifically, we let
x = αs + ξ(r) and y = βs + η(r). The functions ξ and η can be chosen in a number of different ways, and our choice is based
on the desire to have a simple orthogonal coordinate system. In particular, we take
x = αs + βr and y = βs − αr. (319)

With this, ∂s = α∂x + β∂y , and so (317) becomes
∂s u0 + u0 = f (αs + βr, βs − αr). (320)

This is easy to solve and the general solution is
Z s
−s
u0 = a0 (r)e + eτ −s f (ατ + βr, βτ − αr)dτ, (321)
where a0 (r) is arbitrary. Since we have not yet determined exactly what portion of the boundary condition, if any, the outer
solution should satisfy, we are not yet in position to convert back to x, y coordinates.
Before investigating the boundary layer, it is instructive to consider the change of coordinates given in (319). Fixing r
and letting s increase, one obtains the directed lines shown in Fig.37. Thus starting at boundary point P , as the variable
s increases one crosses the domain Ω and arrives at boundary point P∗ In terms of our first-term approximation, we need
to know what value u0 starts with at P and what value it has when the point P∗ is reached. However, we only have one
integration constant in (321), so there is going to be a boundary layer at either P or at P∗ . As we will see later, for our
problem the outer solution will satisfy the boundary condition at P∗ and there will be a boundary layer of width ε at P (see
Fig.37). Therefore, the set ∂Ωo in (318) consists of those points of ∂Ω where the characteristic curves leave Ω as s increases
(like P∗ ). Based on these comments, we will assume that the points of the boundary can be separated into three disjoint
sets: the points where the characteristics leave Ω (this is ∂Ωo ), the points where they enter Ω (this will be designated as
∂Ωb ), and the tangency points (this is ∂Ωt ). For the domain shown in Fig.37 there are two tangency points )i.e., TL and TR ).
These points can cause real headaches when constructing an asymptotic approximation and will be left until the end. Also,
to keep the presentation simple, we will assume that the situation is as pictured in Fig.37 (i.e., the domain is convex so that
any straight line joining any two points on the boundary is wholly inside the boundary). This means that if a characteristic
curve enters the domain, then the only other time it intersects the the boundary is when it leaves the domain. Moreover,
∂Ωo and ∂Ωb are assumed to be smooth curves.
Now that we know (or think we know) where the boundary layer is located, we are in a position to complete the
specification of the outer solution given in (321). To do this, assume ∂Ωo can be described as s = ho (r) for r0 < r < r1 . In
this case, using (318), the outer solution in (321) becomes
75
Z ho (r)
ho −s
u0 = g(αho + βr, βho − αr)e − eτ −s f (ατ + βr, βτ − αr)dτ. (322)
s
This solution does not apply along ∂Ωb or on ∂Ωt . To be more specific, suppose the curve s = hb (r) for r0 < r < r1 describes
∂Ωb . In this case, (322) holds for r0 < r < r1 and hb (r) < s ≤ ho (r).
Boundary-Layer Expansion To find out what goes on in the boundary layer, we introduce the boundary coordinate (see
Fig.38)
s − hb (r)
s= . (323)
ε
This coordinate, by necessity, depends on r since the boundary depends on r. This complicates the calculations in making
the change of variables from (r, s) to (r, s). One finds using the chain rule that the derivatives transform as follows:
1 h′b
∂s → ∂s , ∂r → − ∂s + ∂r ,
ε ε
and
1 2 (h′b )2 2 h′′b 2h′b

∂s2 → ∂s , ∂ 2
r → ∂s − ∂s − ∂s ∂r + ∂r2 .
ε2 ε2 ε ε
To convert the original problem into the boundary-layer coordinates, we must first change from x, y to r, s coordinates
in (313). This is relatively easy since from (319) one finds that
1 1
∂x = (α∂s + β∂r ) and ∂y = (β∂s − α∂r ) ,
γ γ
where γ = α2 + β 2 . With this, (313) becomes

ε ∂s2 + ∂r2 u + γ∂s u + γu = γf. (324)
Now, letting U (r, s) denote the solution in the boundary layer and then substituting the boundary layer coordinates into
(324) yields

µ∂s2 + γ∂s + O(ε) U = εγf, (325)
where
2
µ(r) = 1 + (h′b ) . (326)
In keeping with our usual assumptions, hb (r) is taken to be a smooth function. However, note that at the tangency
points shown in Fig.37, h′b (r) = ∞. For this reason, these points will have to be dealt with separately after we finish with
the boundary layer.
We are now in position to expand the solution in terms of the usual power series expansion, and so let
U (r, s) ∼ U0 (r, s) + . . . . (327)

Substituting this into (325) yields the equation
µ∂s2 U0 + γ∂s U0 = 0.
The general solution of this is
U0 (r, s) = A(r) + B(r)e−γs/µ , (328)

where A(r) and B(r) are arbitrary. Now, from the boundary and matching conditions, we must have U0 (r, 0) = g and
U0 (r, ∞) = u0 (r, hb ). Imposing these on our solution in (328) yields
U0 (r, s) = u0 (r, hb ) + [g (αhb + βr, βhb − αr) − u0 (r, hb )] e−γs/µ , (329)

where Z ho
ho −hb
u0 (r, hb ) = g(αho + βr, βho − αr)e − eτ −hb (r) f (ατ + βr, βτ − αr)dτ.
hb
76
Ο(ε) u
s
hb(r) ho(r)
Figure 38: The solution as a function of the characteristic coordinate s as defined in (319). There is a boundary layer of
width O(ε) at the left end of the interval. Note that in connection with the previous figure hb (r) and ho (r) correspond to P
and P∗ .
Composite Expansion It is not difficult to put together a composite expansion that will give us a first term approximation
of the solution in the outer- and boundary-layer regions. Adding (322) and (329) together and then subtracting their common
part yields
Z ho (r)
u ∼ g(αho + βr, βho − αr)eho −s − eτ −s f (ατ + βr, βτ − αr)dτ (330)
s
+ [g(αhb + βr, βhb − αr) − u0 (r, hb ] e−γs/µ .
The approximation holds for r0 < r < r1 and hb (r) ≤ s ≤ ho (r). This result may not be pretty, but it does give us
a first-term approximation of the solution in the boundary layer and outer domain. What is interesting is that we have
been able to patch together the solution of an elliptic problem using solutions to hyperbolic and parabolic problems. Those
who may be interested in the theoretical foundation of the approximations constructed here are referred to Levinson (1950),
Eckhaus and de Jager (1966), and Il’in (1992).
Given the effort invested in deriving the composite expansion in (330), it is natural to wonder just how well it does in
approximating the solution of the partial differential equation. To demonstrate its accuracy, suppose f = 1, g = 0, and Ω is
the unit disk. The resulting approximation obtained from (330) and the exact solution given in (315) are shown in Fig. 39
(also see Fig). Based on this graph, it seems that we have done reasonably well with our approximation. However, we are
not finished, as (330) does not hold in the immediate vicinity of the tangency points TR and TL , which are shown in figure
37.
Parabolic boundary layer To complete the construction of a first-term approximation of the solution of (313) it remains
to find out what happens near the tangency points shown in Fig.37. We will concentrate on TL , and to do this we let (r0 , s0 )
be its coordinates in the r, s system. Also, suppose the smooth curve r = q(s) describes the boundary ∂Ω in this region. In
this case r0 = q(s0 ) and q ′ (s0 ) = 0. It will be assumed here that q ′′ (s0 ) 6= 0. The boundary-layer coordinates are now
r − q(s) s − s0
re = and se = . (331)
εα εβ
The transformation formulas for the derivatives are similar to those derived earlier for the boundary layer along ∂Ωb , so
they will not be given. The result is that (324) takes the form

ε ε−2α ∂re2 + ε−2β ∂es2 + . . . Ue + γ ε−β ∂es − εβ−α seq0′′ ∂re + . . . Ue + γU
e = γf.
Here Ue (e
r , se) is the solution in this region, and we have used the Taylor series expansion q ′ (s0 + εβ se) ∼ εβ seq0′′ to obtain the
above result. There are at least two balances that need to be considered. One is α = β = 1 and the other is 2β = α = 2/3.
The latter is the one of interest, and assuming
e (e
U e0 (e
r , se) ∼ U r , se) + . . . (332)
one obtains the equation
77
0
-1
u(x,y) -2
Exact
Asymptotic
-3
-1.0 -0.5 0 0.5 1.0
ρ - axis
Figure 39: Comparison between the exact solution (315) and the composite expansion given in (330) with ε = 0.05. The
values of each function are given along the line x = ρ cos(π/4), y = ρ sin(π/4), −1 ≤ ρ ≤ 1.

∂re2 − γe
sq0′′ ∂re + γ∂es Ue0 = 0, (333)
for 0 < re < ∞ and −∞ < se < ∞. This is a parabolic equation, and for this reason this region is referred to as a parabolic
boundary layer. The solution is required to match to the solutions in the adjacent regions and should satisfy the boundary
condition at re = 0. The details of this calculation will not be given here but can be found in van Harten (1976). The papers
by Cook and Ludford (1971, 1973) should also be consulted, as they have an extensive analysis of such parabolic layers and
how they appear in problems where the domain has a corner. The theory necessary to establish the uniform validity of the
expansions when corners are present in the boundary, and a historical survey of this problem, can be found in Shih and
Kellogg (1987).
5.9.3 Leibniz’ rule for differentiating under the integral sign

The idea is to consider:
Z β(t)
f (x, t)dx = I(t, α, β) (334)
α(t)
ˆ and use the chain rule:

and because α = α(t), β = β(t) we can consider I = I(t)
dIˆ ∂I ∂I dα ∂I dβ
= + + .
dt ∂t ∂α dt ∂β dt
Recalling the fundamental theorem of calculus
Z x
d
f (s)ds = f (x)
dx
we note that:
∂I ∂I
= −f (α(t), t); = f (β(t), t)
∂α ∂β
and so
dIˆ ∂I dα dβ
= − f (α, t) + f (β, t) .
dt ∂t dt dt
Comment Leibniz rule holds provided α(t), β(t) and f (x, t) are sufficiently “well behaved”. More precisely: let α(t) and
β(t) be differentiable for t ∈ [t1 , t2 ]. Let A, B be finite constants such that A ≤ α(t) ≤ B and A ≤ β(t) ≤ B on t1 ≤ t ≤ t2 .
If f (x, t) and ft (x, t) are continuous on the rectangle A ≤ x ≤ B, t1 ≤ t ≤ t2 , then the Leibniz rule holds for each t ∈ [t1 , t2 ].
78
5.9.4 Travelling periodic waves and exponential solutions
In linear partial differential equations, for example in one space and time variable, one often has occasion to deal with
solutions of the form:
u(x, t) = A cos(kx − ωt) (335)

referred to as a plane wave or wave train (it is a travelling periodic wave). Here A is the amplitude of the wave, k is the
wave number, and ω is the frequency. The wave number k is a measure of the number of spatial oscillations (per 2π units
of distance) observed at a fixed time, and the (circular) frequency ω is a measure of the number of oscillations in time (per
2π units of time) observed at a fixed spatial location. Some authors use the angular frequency f where ω/2π = f and f is
the number of oscillations per unit of time. The number λ = 2π/k is the wavelength, and P = 2π/ω is the time period. The
wavelength measures the distances between successive crests, and the period measures the time for an observer located at a
fixed position x to see a repeat pattern. The expression may also be written
ω
u(x, t) = A cos k(x − t) (336)
k
so it also represents a travelling wave (f (x − ct)) moving to the right at speed c = ω/k, called the phase velocity. This is the
ω 2π
speed that one would have to move at to remain on the crest of the wave. This can also be written as c = 2π k = f λ. For
calculations the complex form is often preferred:
u(x, t) = A exp[i(kx − ωt)] (337)

because differentiations with the exponential function are simpler. Here A is a complex constant (A = Ar + iAi ). After the
calculations have been completed, one can use Euler’s formula exp(iθ) = cos θ + i sin θ and then take real or imaginary parts
to recover a real solution. (The real and imaginary parts of the complex solution to a linear differential equation are both
real solutions of the equation).
In ordinary differential equations, one often meets forcing (oscillatory in time) terms like A cos(ωt + φ) where A is the
amplitude, ω is the (circular) frequency of the oscillation as defined above i.e., it is the number of oscillations per 2π units
of time and φ is referred to as the phase shift. (ωt + φ is often referred to as being the phase).
6 Multiple Scales
6.1 Introduction
When the domain one is dealing with is of infinite extent, (typically time is the variable one is dealing with and the problems
involved are usually initial value problems), a type of singular behaviour different from boundary layer behaviour can
occur involving the occurrence of multiple (time) scales. When one uses matched asymptotic expansions, the solution is
constructed in different regions that are then patched together to form a composite expansion. The method of multiple
scales differs from this approach in that it essentially starts with a generalized version of a composite expansion. In doing
this, one introduces coordinates for each region (or layer); these new variables are considered to be independent of one
another. A consequence of this is that what may start out as a ordinary differential equation is transformed into a partial
differential equation. Exactly why this helps to solve the problem, rather than make it harder, will be discussed as the
method is developed in this chapter.
The history of multiple scales is more difficult to delineate than, say, boundary-layer theory. This is because the method
is so general that many apparently unrelated approximation procedures are special cases of it. One might argue that the
procedure got its start in the first half of the 19th century. For example, Stokes (1843) used a type of coordinate expansion
in his calculations of fluid flow around an elliptic cylinder. Most of these early efforts were limited, and it was not until
the latter half of the 19th century that Poincaré (1882), made more extensive use of the ideas underlying multiple scales in
his investigations into the periodic motion of planets. He found that the approximation obtained from a regular expansion
accurately described the motion for only a few revolutions of the planet, after which the approximation was way off. The error
was due, in part, to the contributions of the second term of the expansion. He referred to this difficulty as the presence of a
secular term. To remedy the situation, he expanded the independent variable with the intention of making the approximation
uniformly valid by removing the secular term. This idea is also at the heart of the modern version of the method. What
Poincaré was missing was the introduction of multiple independent variables based on the expansion parameter. This step
came much later; the most influential work in this regard has been by Kuzmak (1959) and Cole and Kevorkian (1963).
6.2 A note on resonance

Consider the linear oscillator (simple harmonic motion) equation with a periodic forcing term
y ′′ + k 2 y = cos ωt. (338)

The solution of the homogeneous equation (i.e., the complementary function) is
79
y = A cos kt + B sin kt
and this suggests that the natural frequency of the oscillator is k i.e., the oscillator would naturally vibrate at that frequency
in the absence of other effects. The forcing term in (338) has a frequency ω; if ω = k it is natural to ask what happens. In
this case the RHS forcing term is forcing the equation at its own natural frequency (this is what a diver on a diving board
does) and we might expect that this would give to rise to an oscillation of ever increasing amplitude. This phenomenon is
known as resonance. (The forcing term is continually reinforcing the natural oscillations of the system).
To illustrate the idea mathematically let us set ω = 1. If we now take k = 1, then we recall that the general solution is
y = yc + yp where yc is the complementary function and yp is any particular solution of the full equation. It is easy to show
that the general solution (verify with Maple) of (338) is
1 1 1
y = A cos t + B sin t + cos t + t sin t = C cos t + B sin t + t sin t.
2 2 2
(To find yp (t) without Maple, try the substitution yp (t) = At cos t + Bt sin t).The constants B, C would be determined by
the (two) boundary conditions of the particular problem.
The most interesting term is the last one which breaks the periodicity of the solution. A little reflection suggests that as
t grows, this term will carry on increasing in absolute value indefinitely. It is referred to as a resonant or secular term
and it corresponds to the case where the forcing term has the same frequency as the unforced oscillator so the solution grows
and grows. This will always happen if the frequency of the forcing term (ω) is the same as the natural frequency (k) of the
oscillator. More generally, resonance will occur if the RHS forcing term contains a term proportional to any of
the homogeneous solutions of (338).
For example consider the case k = 1, ω = 2. Maple shows that the solution of the equation is now
1
y = A cos t + B sin t − cos 2t
3
and there is no secular term.
As a final caution consider the problem
y ′′ + y = cos3 t
Here the RHS is apparently not proportional to cos t or sin t. Yet Maple shows that the solution is
1 3 3
cos3 t + cos t + t sin t
y = A cos t + B sin t −
8 8 8
and a secular term is occurring. Why? If we recall the trigonometric result
cos 3t = 4 cos3 t − 3 cos t

we see that cos3 t actually contains a term proportional to cos t. This causes the resonance. As we will see shortly, it is
necessary to understand these ideas in order to use the multiple scales technique.
Fredholm A more precise mathematical treatment (see appendix 10.5.2) requires use of the Fredholm alternative theorem
which in the context of periodic solutions of the equation:
utt + u = f (t),
states that periodic solutions (the statement that solutions are periodic is in effect a boundary condition) can only exist if
the inhomogeneous term on the RHS is orthogonal to the null space of the (self adjoint) differential operator L[u] = utt + u.
The null space is just the set of all solutions of the homogenous equation utt + u = 0. In the present instance this means
that there can only be periodic solutions (i.e., non-secular solutions) provided:
Z T Z T
f (t) sin t dt = 0, f (t) cos t dt = 0
0 0
where T = 2π is the period in this example. Recalling the familiar orthogonality results from Fourier series theory, if we
express f (t) in a Fourier series, we see that both of these integrals will be zero unless f (t) contains a term either sin t or cos t.
As in the last chapter, we will introduce the ideas underlying the method by going through a relatively simple example.
80
6.3 Example
Consider the problem:
y ′′ (t) + y(t) = εy(t), for 0 < t, (339)

with the following initial conditions:
y(0) = 0 and y ′ (0) = 1. (340)

It is easy to show that a 2 term regular perturbation approach yields a solution which demonstrates secular growth.
6.3.1 Regular perturbation approach

Let us try
y(t) ∼ y0 (t) + εy1 (t) + . . . ,

and the equation becomes:
y0′′ + εy1′′ + y0 + εy1 = εy0 + O(ε2 )

while the boundary conditions become:
y0 (0) + εy1 (0) = 0 and y0′ (0) + εy1′ (0) = 1

In the usual way we get a sequence of problems:
ε0 : y0′′ + y0 = 0, y0 (0) = 0, y0′ (0) = 1

and this has the solution
y0 = sin t.
At the next order we have:
ε1 : y1′′ + y1 = y0 (= sin t), y1 (0) = 0, y1′ (0) = 0.

We notice immediately that there will be resonant terms because of the RHS. Specifically, the solution of this problem
consists of a complementary function and a particular solution. In fact:
1
y1 = A cos t + B sin t − t cos t
2
and applying the boundary conditions we find that A = 0, B = 1/2. The regular perturbation solution is thus
1 1
y ∼ sin t + ε( sin t − t cos t)
2 2
We notice the presence of secular terms. In the ε1 term, observe that sin t and cos t are O(1) functions for all values of
t so that ε sin t remains O(ε) for all t. However 2ε t cos t is more complicated. As t grows cos t remains O(1) but t does not.
And in fact the solution will lose its well orderedness when 2ε t cos t becomes as large as the ε0 term sin t. Specifically, the
second term becomes as large as the first when εt = O(sin t) = O(1) i.e., when t = O(1/ε). So this solution is only valid
when t ≪ 1/ε. It fails when t = O(1/ε).
Let us try a multiple scales approach to see if we can improve on the regular perturbation approach. In particular multiple
scales will allow us to find an approximation which is asymptotically valid on a longer time interval.
6.3.2 Multiple scales approach

This problem has an oscillatory component of the solution
that occurs on a time scale that it O(1) and also has a slow variation
in the solution that takes place on a time scale O 1ε . To incorporate two time scales into the problem, we introduce the
variables
t1 = t (fast time) (341)

and
t2 = εt (slow time) (342)
81
These two time scales will be treated as independent. In constructing our solution, the slowly changing features will be
combined into factors which are functions of t2 , while the rapidly changing features will be combined into factors which are
a function of t1 . Thus we transform y(t) → v(t1 , t2 ). It is good practice to use a different letter for the dependent variable
to emphasise the different dependencies though we will not always be so careful. Note that t1 and t2 are not actually
independent but v is certainly a solution of some equation in which t1 and t2 figure as independent variables. Note also
that there are infinitely many functions v(t1 , t2 ) which equal y(t) along the line t2 = εt1 in (t1 , t2 ) space. So we now have
freedom to prescribe extra additional conditions. With the unwelcome appearance of secular terms, it is natural to think of
conditions chosen such that particular terms (e.g. secular terms) are omitted when we construct an approximation.
Thus using the chain rule the original time derivative transforms as follows:
dy dt1 ∂v dt2 ∂v ∂v ∂v
→ + = +ε
dt dt ∂t1 dt ∂t2 ∂t1 ∂t2
or in operator notation
d ∂ ∂v d2 ∂2 ∂2 ∂2
= +ε 2
= 2 + 2ε + ε2 2
dt ∂t1 ∂t2 dt ∂t1 ∂t1 ∂t2 ∂t2
Substituting this into (362) yields

∂t21 + 2ε∂t1 ∂t2 + ε2 ∂t22 v + v = εv, (343)
where
v = 0 and (∂t1 + ε∂t2 ) v = 1, for t1 = t2 = 0. (344)

∂
To simplify the notation, we use the symbol ∂t1 in place of (and similarly for ∂t2 ). The benefits of introducing these
∂t1
two time variables are not yet apparent. In fact, one might argue that we have made the problem harder since the original
ordinary differential equation has been turned into a partial differential equation. The reasons for doing this will become
evident in the discussion that follows. It should be pointed out that the solution of (343) is not unique and that we need to
impose more conditions for uniqueness on the solution. This freedom will enable us to prevent secular terms from appearing
in the expansion (at least over the time scales we are using).
We use a power series expansion of the form
v ∼ v0 (t1 , t2 ) + εv1 (t1 , t2 ) + · · · . (345)

with the requirement that the series be asymptotic for t2 = Os (1). Substituting this into (371) yields the following

∂t21 + 2ε∂t1 ∂t2 + · · · (v0 + εv1 + · · · ) + v0 + εv1 = ε(v0 + εv1 ). (346)
In writing equation (346), we made an effort to include only the terms that might contribute to the first two terms of the
expansion. From this and (344), we obtain the first problem to solve:

ε0 ∂t21 + 1 v0 = 0, v0 = 0, ∂t1 v0 = 1 at t1 = t2 = 0.
The general solution of this pseudo PDE problem is
v0 = a0 (t2 ) sin(t1 ) + b0 (t2 ) cos(t1 ), (347)

where
a0 (0) = 1 and b0 (0) = 0. (348)

For the moment, the coefficients a0 (t2 ) and b0 (t2 ) are arbitrary functions of t2 except that they are required to satisfy the
initial conditions given in (348).
With the initial conditions in (344), we get the following problem:

ε1 ∂t21 + 1 v1 = v0 − 2∂t1 ∂t2 v0 , v1 = 0 and ∂t1 v1 = −∂t2 v0 when t1 = t2 = 0.
From (347), the differential equation for v1 is

∂t21 + 1 v1 = (2b′0 + a0 ) sin(t1 ) − (2a′0 − b0 ) cos(t1 ). (349)
The general solution of this problem is found to be
1 1
v1 = a1 (t2 ) sin(t1 ) + b1 (t2 ) cos(t1 ) − (2b′0 + a0 ) t1 cos(t1 ) − (2a′0 − b0 ) t1 sin(t1 ), (350)
2 2
where
82
a1 (0) = b1 (0) = 0. (351)
In (350) we have the possibility of secular terms in the expansion. However, the functions a0 and b0 can be chosen to prevent
this. In fact, even if we did not have this solution we see from (349) that to exclude resonating terms we must set:
2b′0 + a0 = 0; 2a′0 − b0 = 0 (352)

and this linear constant coefficient system must be solved subject to boundary conditions (348). From the first equation we
see that:
a0 = −2b′0 (353)
and the second equation now gives
4b′′0 + b0 = 0 (354)
so that
t2 t2
b0 = e cos + f sin (355)
2 2
From (352), we see that
1
2b′0 (0) = −a0 (0) ⇒ b′0 (0) = − ; b0 (0) = 0 (356)
2
and so e = 0, f = −1 and
t2
b0 = − sin (357)
2
From (359) we see that
t2
2b′0 + a0 = 0 ⇒ a0 = cos (358)
2
It is worth emphasising that it is not actually necessary to solve (349). Secular terms are avoided by simply removing the
solutions of the associated homogenous equation in (349). This means we require the coefficients of sin(t1 ) and cos(t1 ) in
(349) to be zero, which leads to the same conclusions given in (352).
Putting our results together, we have found that in terms of the original variables
t2 t2 εt εt
v ∼ cos sin t1 − sin cos t1 =⇒ y(t) = cos sin t − sin cos t. (359)
2 2 2 2
This is a one term approximation that is valid up to at least t2 = O(1) i.e., t = O(1/ε) and in fact the solution appears to
be valid for all t. This certainly improves upon our regular perturbation which required t ≪ O(1/ε).
Note This (second order linear homogeneous constant coefficient) equation has the general exact solution (exercise):
y(ε, t) = a cos((1 − ε)1/2 t + α) (360)

where a, α are arbitrary constants to be determined from the boundary conditions.
6.4 Derivation of oscillator equation with weak damping

Suppose we are dealing with a mass-spring-dashpot system as in fig. 40 where x(t) denotes the displacement of the mass m0
upwards from the position it is in when the spring is unextended. When the spring is extended, the forces on the mass are
the spring force Fs , the damping force exerted by the dashpot Fd . If the spring is linear then the force which it exerts is
′
Fs = −kx(t) where k is the spring constant (Hooke’s law). If the damping is also linear then Fd = −cx (t) i.e., the damping
is proportional to the velocity of the mass only. If the mass is in equilibrium (i.e., not moving) then x(t) = xs (a constant),
the damping force will be zero and the gravity and spring forces balance one another. Hence
m0 g
Fs = m0 g =⇒ −kxs − m0 g = 0 =⇒ xs = − .
k
So in the coordinate system chosen, the equilibrium position corresponds to a negative displacement (downwards). Now
suppose the system starts from its equilibrium position and is disturbed with a known initial velocity initial velocity v0 . The
corresponding initial value problem for the displacement of the mass is
′′ ′ ′
m0 x = −kx(t) − m0 g − cx (t); x(t = 0) = xs ; x (t = 0) = v0.
83
x(t)
spring(k)
mass m0
damper(c)
Figure 40: A mass bobbing on the end of a spring (with damping). In practice, even without the damping element, a system
such as this will experience damping due to friction effects. This is an example of a damped linear oscillator (damped simple
harmonic motion).
For convenience we shift the origin of the vertical coordinate to the equilibrium position. Thus let y(t) ≡ x(t) − xs and
′ ′′ ′′
we see that y ′ (t) ≡ x (t) and y (t) = x (t) and the problem becomes:
′′ ′ ′′ ′ ′
m0 y = −k(xs + y(t)) − m0 g − cy (t) =⇒ m0 y = −ky(t) − cy (t); y(0) = 0; y (t = 0) = v0 .
This is the formulated initial value problem in dimensional form. We now non-dimensionalise using the procedure described
in the first chapter. We first list the variables and the parameters:
variables units
dependent variable y m
parameters
damping constant c Nsm−1
spring constant k Nm−1 =Kgs−2
initial velocity v0 ms−1
mass m0 Kg
Note that gravity does not actually appear in the formulated problem. We non-dimensionalise as we would non-dimensionalise
the undamped system (as if c = 0) i.e., we non-dimensionalise so as to balance the first two terms in the equation. Specifically,
let us introduce dimensionless variables y, t :
y = yL; t = tT
where T, L are time and length scales to be chosen via the basic balances in the equation. The dimensionless problem is:
m0 L d2 y cL dy L dy
2 2 = −kLy − , y(t) = 0, = v0 .
T dt T dt T dt
We now choose to balance the first two terms in the equation and this requires:
r
m0 L m0
= kL =⇒ T ≡
T2 k
p m0
which fixes the time scale. (Verify that k is indeed a time scale). Consider now the inhomogeneous boundary condition
and require that both terms balance so that:
r
L m0
= v0 =⇒ L = v0 T = v0 .
T k
Recall that F = ma so N=Kgms−2 and verify that this does define a lengthscale.
84
Thus we introduce:
r r
m0 m0
y = yv0 ;t = t
k k
and the problem becomes
r r r r
m0 k d2 y m0 m0 k dy
m0 v0 = −kv0 y(t) − cv0
k m0 dt2 k k m0 dt
which simplifies to:
d2 y c dy dy(t = 0)
2 = −y(t) − √ ; y(t = 0) = 0; = 1.
dt km0 dt dt
c
N.B. Do not forget to transform the initial conditions. If the dimensionless parameter √km is small (all these
0
c
parameters are known quantities), then we are dealing with a system with weak damping and we define √km ≡ ε ≪ 1 to
0
get the system in its final dimensionless form:
d2 y dy dy(t = 0)
2 = −y(t) − ε ; y(t = 0) = 0; = 1. (361)
dt dt dt
6.5 Example
The problem to be considered is the oscillator equation (361) where we drop the bars on the variables. We wish to find the
y(t) that satisfies
y ′′ + εy ′ + y = 0, for 0 < t, (362)

where
y(0) = 0 and y ′ (0) = 1. (363)

This equation corresponds to an oscillator (i.e., a mass-spring-dashpot) with weak damping, where the time variable has
been scaled by the period of the undamped system. This is the classical example used to illustrate multiple scales. Our
approach to constructing the approximation comes from the work of Cole and Kevorkian (1963) and Reiss (1971).
6.5.1 Regular Expansion

It appears that this is not a layer type of problem since ε is not multiplying the highest derivative in (362). It seems,
therefore, that the solution has a regular power series expansion. So, assuming
y(t) ∼ y0 (t) + εy1 (t) + . . . , (364)

one finds from (362) that
1
y(t) ∼ sin(t) − εt sin(t). (365)
2
In comparison, the exact solution of the problem is
1 p
y(t) = p e−εt/2 sin t 1 − ε2 /4 . (366)
1 − ε2 /4
These functions are plotted in Fig.41. It is apparent that we have not been particularly successful in finding a very
accurate approximation of the solution for 0 ≤ t ≤ ∞. This is because the second term in the expansion is as large as the
first term once εt ≈ 1. For this reason, the second term is said to be a secular term, and for the expansion to be well ordered
we must have εt ≪ 1. The problem with our approximation is that the solution (366) decays but the first term in (365) does
not. The second term tries to compensate for this, and in the process it eventually becomes as large as the first term.
Another way to look at what is happening in the problem is to consider the energy. By multiplying (362) by the velocity
y ′ , integrating and then using the initial conditions (363), one obtains the equation.
Z t
2
H(y, y ′ ) = H(0, 1) − ε [y ′ (τ )] dτ, (367)
0
where
1 ′ 2 1 2
H(y, y ′ ) ≡ (y ) + y . (368)
2 2
85
2 Exact Solution
1 y0(t) + εy1(t)
Solution
0
-1
-2
0 25 50 75
t - axis
Figure 41: Regular perturbation expansion and exact solution (366) with ε = 0.1. There is reasonable agreement up to about
t ≈ 10 but they differ completely for larger t.
The function H is the Hamiltonian and consists of the sum of the kinetic and potential energies for the system. Thus, (367)
shows that the energy equals the amount included at the start minus whatever is lost due to damping. To relate this to the
multiple-scale expansion, recall that a system is said to be conservative, or energy is conserved, if H is independent of t. On
d
the other hand, if dt H < 0, so that energy is lost, then the system is dissipative. Using these ideas, we have that (362) is
dissipative but the problem for y0 (t) is conservative. Therefore, the first term in the approximation in (365) is a conservative
approximation to a dissipative problem. This is one of many situations where the multiple-scales method is generally needed.
6.5.2 Multiple-Scale Expansion

This problem has an oscillatory component of the solution
that occurs on a time scale that it O(1) and also has a slow variation
in the solution that takes place on a time scale O 1ε . To incorporate two time scales into the problem, we introduce the
variables
t1 = t (fast time of oscillation) (369)

and
t2 = εα t (slow time of the amplitude drift). (370)

Based on the earlier discussion of the exact solution, it is expected that α = 1. However, we will derive this later when
constructing the expansion of the solution. These two time scales will be treated as independent. In constructing our solution,
the slowly changing features will be combined into factors which are functions of t2 , while the rapidly changing features will
be combined into factors which are a function of t1 . Thus we transform y(t) → v(t1 , t2 ). It is good practice to use a different
letter for the dependent variable to emphasise the different dependencies though we will not always be so careful. Note that
t1 and t2 are not actually independent but u is certainly a solution of some equation in which t1 and t2 figure as independent
variables.
Thus using the chain rule the original time derivative transforms as follows:
dy dt1 ∂v dt2 ∂v ∂v ∂v
= + = +ε
dt dt ∂t1 dt ∂t2 ∂t1 ∂t2
or in operator notation
d ∂ ∂v d2 ∂2 ∂2 ∂2
= +ε 2
= 2 + 2ε + ε2 2
dt ∂t1 ∂t2 dt ∂t1 ∂t1 ∂t2 ∂t2

∂t21 + 2εα ∂t1 ∂t2 + ε2α ∂t22 v + ε (∂t1 + εα ∂t2 ) v + v = 0, (371)
where
86
v = 0 and (∂t1 + εα ∂t2 ) v = 1, for t1 = t2 = 0. (372)
To simplify the notation, we use the symbol ∂t1 in place of ∂t∂1 (and similarly for ∂t2 ). The benefits of introducing these
two time variables are not yet apparent. In fact, one might argue that we have made the problem harder since the original
ordinary differential equation has been turned into a partial differential equation. The reasons for doing this will become
evident in the discussion that follows. It should be pointed out that the solution of (371) is not unique and that we need to
impose more conditions for uniqueness on the solution. This freedom will enable us to prevent secular terms from appearing
in the expansion (at least over the time scales we are using).
We use a power series expansion of the form
v ∼ v0 (t1 , t2 ) + εv1 (t1 , t2 ) + · · · . (373)

with the requirement that the series be asymptotic for t2 = Os (1). Substituting this into (371) yields the following

∂t21 + 2εα ∂t1 ∂t2 + · · · (v0 + εv1 + · · · ) + ε (∂t1 + · · · ) (v0 + · · · ) + v0 + εv1 + · · · = 0. (374)
expansion. From this and (372), we obtain the first problem to solve:

ε0 ∂t21 + 1 v0 = 0, v0 = 0, ∂t1 v0 = 1 at t1 = t2 = 0.
The general solution of this pseudo PDE problem is
v0 = a0 (t2 ) sin(t1 ) + b0 (t2 ) cos(t1 ), (375)

where
a0 (0) = 1 and b0 (0) = 0. (376)

For the moment, the coefficients a0 (t2 ) and b0 (t2 ) are arbitrary functions of t2 except that they are required to satisfy
the initial conditions given in (376).
To go any farther, we need to determine α in the slow time scale (370). To do this, note that in the O(ε) equation
obtained from (374) we are going to have the term ∂t1 v0 . It is not hard to check that this will cause a secular term to appear
in the expansion. However, this can be prevented if we make use of the dependence of v0 on the time scale t2 . In (374) the
only term available to us that can be used to prevent a secular term is 2εα ∂t1 ∂t2 v0 . Balancing this with ε∂t1 v0 gives us α = 1.
With this, and the initial conditions in (372), we get the following problem:

ε1 ∂t21 + 1 v1 = −2∂t1 ∂t2 v0 − ∂t1 v0 , v1 = 0 and ∂t1 v1 = −∂t2 v0 when t1 = t2 = 0.
Using (375), the (pseudo PDE) differential equation for v1 is

∂t21 + 1 v1 = (2b′0 + b0 ) sin(t1 ) − (2a′0 + a0 ) cos(t1 ). (377)
The general solution of this problem (to find the particular solution for this inhomogeneous problem the substitution
vp (t1, t2 ) = A(t2 )t1 cos t1 + B(t2 )t1 sin t1 can be used) is found to be
1 1
v1 = a1 (t2 ) sin(t1 ) + b1 (t2 ) cos(t1 ) − (2b′0 + b0 ) t1 cos(t1 ) − (2a′0 + a0 ) t1 sin(t1 ), (378)
2 2
where
a1 (0) = b1 (0) = 0. (379)

In (378) we have the possibility of secular terms in the expansion. However, the functions a0 and b0 can be chosen to prevent
this. From (378) and (376), we find that
2b′0 + b0 = 0 b (t ) = β0 e−t2 /2 b0 (t2 ) = 0

′ ⇒ 0 2 ⇒ . (380)
2a0 + a0 = 0 a0 (t2 ) = α0 e−t2 /2 a0 (t2 ) = e−t2 /2
In the last step, we impose the initial conditions given in (376).
It is worth pointing out that it is not actually necessary to solve (377). Secular terms are avoided by simply removing
the solutions of the associated homogenous equation in (377). This means we require the coefficients of sin(t1 ) and cos(t1 )
in (377) to be zero, which leads to the same conclusions given in (380).
Putting our results together, we have found that (in terms of the original variables)
y ∼ e−εt/2 sin(t). (381)
87
1.0
Exact Solution
0.5 y0(t1,t2)
Solution
0
-0.5
-1.0
0 20 40 60 80 100
t - axis
Figure 42: Multiple scale approximation (381) and the exact solution. The agreement is excellent. ε = 0.1.
This is a first-term approximation that is valid up to at least t2 = O(1) i.e., t = O(1/ε). This improves upon our regular
perturbation (365) which required t ≪ O(1/ε). This approximation and the exact solution are shown in Fig. 42, from
which it is clear that the multiple-scale result is quite good. In fact, one can prove that this is a uniformly valid asymptotic
approximation of the solution for 0 ≤ t ≤ O 1ε (see Exercises).
Note We also note that that an alternative argument (Kevorkian/Cole 1985) is sometimes used to remove the secular
terms. If we consider (378), we notice that the O(ε) second term in the expansion could be written
1 1
εv1 = εa1 (t2 ) sin(t1 ) + εb1 (t2 ) cos(t1 ) − (2b′0 + b0 ) εt1 cos(t1 ) − (2a′0 + a0 ) εt1 sin(t1 ) (382)
2 2
and in particular the last two terms can be rewritten as:
1 1
(2b′0 + b0 ) t2 cos(t1 ) − (2a′0 + a0 ) t2 sin(t1 ) (383)
2 2
using the fact that t2 = εt1 . But now this term is formally O(1) rather than O(ε) as we require the solution to be asymptotic
for t2 = Os (1) and the series is not well ordered. Thus we must choose a0 , b0 according to (380) so that the series remains
well ordered.
6.5.3 Three Time Scales

What if one wants an approximation that holds up to ε2 t = O(1)? In this case, the third time scale t3 = ε2 t needs to be
introduced. The difference now is that the coefficients a0 and b0 in (375) depend on both t2 and t3 . In particular, to remove
secular terms in (378), one obtains
a0 (t2 , t3 ) = α0 (t3 )e−t2 /2 and b0 (t2 , t3 ) = β0 (t3 )e−t2 /2 ,

where α0 (0) = 1 and β0 (0) = 0. To find α0 (t3 ) and β0 (t3 ), it is necessary to solve the O(ε2 ) equation, which is

∂t21 + 1 v2 = −2∂t1 ∂t2 v0 − ∂t1 v0 − ∂t1 v1 − ∂t22 + 2∂t1 ∂t3 v0 .
Substituting the solutions for v0 and v1 given in (375) and (378) shows that to prevent secular terms in the expansion
we must require that
1
a1 = α1 (t3 )e−t2 /2 + (β0 − 8α′0 ) t2 e−t2 /2 (384)
8
and
1
(α0 + 8β0′ ) t2 e−t2 /2 .
b1 = β1 (t3 )e−t2 /2 − (385)
8
This is a different situation than before, as we have apparently removed the secular terms yet still have not determined α0
and β0 . These coefficients can be found if we remember that our objective is to obtain an accurate asymptotic approximation
- this means minimizing the error. This is discussed in more detail later, but for the moment note that our expansion is
88

1
v ∼ α0 (t3 ) + ε(α1 (t3 ) + (β0 − 8α0 ) t2 ) e−t2 /2 sin(t1 )
′
8

1
+ β0 (t3 ) + ε(β1 (t3 ) − (α0 + 8β0 t2 ) e−t2 /2 cos(t1 ).
′
8
To prevent the O(εt2 ) terms from interfering with the validity of our expansion, we take β0 − 8α′0 = 0 and α0 + 8β0′ = 0.
From this, and the initial conditions (376), it follows that α0 = cos(t3 /8) and β0 = − sin(t3 /8). Therefore, a first-term
approximation valid up to ε2 t = O(1) is
y ∼ e−εt/2 sin ωt, (386)

1 2
where ω = 1 − 8ε .
6.5.4 Discussion and Observations

1. Now that we have gone through the analysis of a problem using multiple scales, it is worth reflecting on what we did.
The approach used to construct the expansion in (381) was to introduce the two time scales t1 = t and t2 = εt. For this
reason, the method is sometimes called two-timing, and t1 is said to be the fast time scale and t2 the slow scale. This
terminology implies that the method is used only with the time variable, which is definitely not the case. In general,
multiple scales are most useful when two disparate scales are operating simultaneously in a problem.
Nevertheless we will see in 6.7 that it is possible to solve some of the boundary-layer problems of the last chapter by
using these same ideas.
The time scales that should be used depend on the problem and the interval over which one wants an approximation
of the solution. The list of possible scales is endless, but some possibilities include
(a) Several time scales: for example, one may need to use the four time scales: t1 = t/ε, t2 = t, t3 = εt, t4 = ε2 t.
(b) More complex dependence on ε: for example t1 = (1+ω1 ε+ω2 ε2 +· · · )t and t2 = εt, where the ωi ’s are determined
while solving the problem (this is called the method of strained coordinates (see end of chapter) or Lindstedt’s
method or the Poincare-Linstedt method ). In the previous problem, one would find that t1 = 1 − ε2 /8 + · · · t.
(c) The correct scaling may not be immediately apparent, and one starts off with something like t1 = εα t and t2 = εβ t,
where α < β.
(d) Nonlinear time dependence: for example, one may have to assume t1 = f (t, ε) and t2 = εt1 , where the function
f (t, ε) is determined from the problem (this type of scaling is used in the next section).
2. Multiple scales are introduced to keep the expansion well ordered. This should not be interpreted to mean that
they are always to remove unbounded terms, although many times this is what happens (as in the preceding example).
In fact, whatever freedom there is in the dependence on t2 is used to minimize the error in the approximation.
We did this in the previous example by trying to reduce the contribution of the next term in the expansion. The
coefficients a0 and b0 in the first-order approximation in (375) were determined by requiring y1 to be as small as we
could make it over the time scales we were using. This minimization principle is, in fact, the condition imposed in the
multiple-scale method.
To illustrate the principle of minimum error, suppose the exact solution is found to be y(t) = 1 + ε2 t + ε4 sin(t). Also,
suppose we use the time scales t1 = t, t2 = εt, and t3 = ε2 t. In this case, we can take y ∼ y0 + εy1 + ε2 y2 + · · · , where
(a) y0 = 1, y1 = 0, y2 = t1 ,
(b) y0 = 1, y1 = t2 , y2 = 0,
(c) y0 = 1 + t3 , y1 = 0, y2 = 0
(d) y0 = 1 + t3 , y1 = −t2 , y2 = t1 .
Each one of these is correct. However, in determining y0 , we want to minimize the value of the error

E0 = max y(t) − y0 (t, εt, ε2 t) ,
where the max is taken over the interval 0 ≤ t ≤ O(1/ε2 ) (i.e., over the time interval 0 ≤ t ≤ T /ε2 where T is fixed).
In the above example, this means y0 = 1 + t3 . Once this result is established, then in determining y1 we want the error

E1 = max y(t) − y0 (t, εt, ε2 t) − εy1 (t, εt, ε2 t)
to be minimized. One finds that y1 = 0, and the other terms can be obtained using a similar procedure. For many
problems, however, it is not an easy matter to actually prove that the error is minimized. There are also questions
associated with what norm to use and the completeness of the time scales. Some of these issues are discussed in Smith
(1985).
89
3. To illustrate another aspect of multiple scales, suppose one has a problem that contains the three time scales t1 =
t, t2 = εt, and t3 = ε2 t. Also suppose the following expansions have been constructed:
(a) y ∼ y0 (t1 , t2 ),
(b) y ∼ y0 (t1 , t2 ) + εy1 (t1 , t2 ),
(c) y ∼ y0 (t1 , t2 , t3 ).
With several different approximations of the same function, it is natural to ask how they compare. Well, the difference
in (i) and (iii) is that (i) is guaranteed to hold up to εt = O(1) while (iii) holds up to ε2 t = O(1). There is no reason
to expect, however, that (iii) is more accurate than (i) over the time interval where they both apply, that is, for
0 ≤ εt ≤ O(1). On the other hand, (ii) should be a more accurate approximation than either (i) or (iii) as ε becomes
small and 0 ≤ εt ≤ O(1). Note that the effort to find (ii) and (iii) is likely to be about the same since the O(ε2 ) problem
must be examined to be able to find y1 as well as remove the secular terms associated with the t3 time scale.
Nonlinear problems
Before we close this section, it is worth commenting on the equation
y ′′ + ω 2 y = 0, (387)
where ω is a positive constant. The general solution of this can be written as
y(t) = a0 sin(ωt) + b0 cos(ωt), (388)

or (with a0 = A cos θ, b0 = −A sin θ)
y(t) = A cos(ωt + θ) = A cos θ cos ωt − A sin θ sin ωt (389)

where A is the amplitude of the oscillation, ω is the frequency and θ is called the phase shift. In many of the examples and
exercises in this chapter, equation (387) must be solved to obtain the first term in the approximation. Indeed, this is exactly
what happened in the earlier example, and it resulted in the solution given in (375). In this example, where the problem
given in (362) is linear, either form of the general solution can be used without difficulty. However, this usually is not the
case when the original problem is nonlinear. It is recommended that whenever using multiple scales on a nonlinear problem
and (387) appears, the solution be written as in (389). With this representation, A and θ are functions of the slower time
variables (just as a0 and b0 are when using (388)). The reason for using the more compact form given in (389) is that it can
greatly simplify certain calculations, such as finding powers of y n . Some formulas that prove useful for such situations are
given in Holmes(1995), p.114.
6.6 Slowly Varying Coefficients

To illustrate what happens when a more complicated time scale must be used, consider the following problem:
y ′′ + k 2 (εt)y = 0, for 0 < t, (390)

where
y(0) = a and y ′ (0) = b. (391)

This equation corresponds to an undamped oscillator with a spring that has a slowly varying restoring force. This is the sort
of situation that could arise, for example, if the spring gradually fatigued or changed periodically in time. We will assume
that k(τ ) is a smooth positive function for τ ≥ 0. It is also assumed that it is not constant.
Before deriving a multiple-scales approximation to the solution of (390), it is worth examining what a solution can look
like. Therefore in Fig. 43 the numerical solution is shown for a particular choice of the function k(εt). It is evident from
this graph that there are (at least) two time scales. There is a fast time scale associated with the rapid oscillations in the
solution, and there is a slower time scale involved with the variation of the amplitude of the oscillation. Our objective is to
determine exactly what these time scales are and how they depend on the coefficient function k(εt).
Based on our success with the weakly damped oscillator in (6.5), suppose we assume the two principal time scales for
this problem are t1 = t, t2 = εt and y(t) → v(t1 , t2 ). Using the same expansion as in (373), one finds that
v0 = a0 (t2 ) sin [k(t2 )t1 ] + b0 (t2 ) cos [k(t2 )t1 ] .

To determine a0 and b0 , we need to solve the O(ε) equation. One finds that this equation is

∂t21 + k 2 v 1 = −2 (a0 k)′ − b0 kk ′ t1 cos [k(t2 )t1 ] + 2 (b0 k)′ + a0 kk ′ t1 sin [k(t2 )t1 ] (392)
90
0.6
0.3
y(t)
0
-0.3
-0.6
0 25 50 75 100 125
t - axis
Figure 43: Dashed curve shows the numerical solution of (390) with k(εt) = 2 + cos(εt), ε = 0.1, a = 0, b = 1. The solid
curve shows the amplitude of the oscillation determined from the multiple scales approximation (404).
Looking at the right-hand side of this equation, one sees that terms like t21 sin [k(t2 )t1 ] and t1 sin [k(t2 )t1 ] are going to
appear in the solution for v1 . Since a0 , b0 and k are independent of t1 , the only way to prevent this is to take a0 = b0 = 0. In
other words, the only solution we are able to find that does not contain secular terms is the zero solution. Clearly, we have
not gotten very far with our simple expansion.
It is necessary to think a little harder on how to choose the time scales, and the first place to begin is with t1 . This time
scale should be a measure of the period of oscillation. In this problem, the period depends on k(εt). (Note for example that
′′
if we consider y = −k 2 y where k is a constant, then the solution is of the form y = A cos kt + B sin kt.) At the moment, it
is unclear exactly how to define this time scale and so we will let
t1 = f (t, ε) (fast) (393)

and
t2 = εt (slow). (394)
The function f (t, ε) is required to satisfy the following conditions:
1. f (t, ε) is nonnegative and increases with t,

2. εt ≪ f as ε ↓ 0 (i.e., t1 is the fast time scale and t2 is the slow scale)
3. f (t, ε) is smooth.
Now, with the multiple-scale transformation and the chain rule, y(t) = v(t1 , t2 ) we have
d df ∂ ∂
= +ε
dt dt ∂t1 ∂t2
and
d2 2 ∂
2
∂ ∂2 2 ∂
2
= f t + f tt + 2εf t + ε .
dt2 ∂t21 ∂t1 ∂t1 ∂t2 ∂t22

ft2 ∂t21 + ftt ∂t1 + 2εft ∂t1 ∂t2 + ε2 ∂t22 v + k 2 (εt)v = 0. (395)
(396)
The oscillation takes place because of the balance between the first and last terms in the equation. Thus we simply take
ft = k(εt). Integrating this yields
91
Z t
f (t, ε) = k(ετ )dτ. (397)
0
It is not hard to show that this choice satisfies the three requirements listed earlier.2 Thus because εk ′ (εt) ≡ εk ′ (t2 ) we have
ftt = εdk/dt2 . So, with the expansion
v ∼ v0 (t1 , t2 ) + εv1 (t1 , t2 ) + · · · , (398)

we obtain the following from (395) and (391):
ε0 (∂t21 + 1)v0 = 0; v0 (0, 0) = a, k(0)∂t1 v0 (0, 0) = b. (399)

The general solution of this problem is
v0 = a0 (t2 ) sin(t1 ) + b0 (t2 ) cos(t1 ), (400)

where a0 (0) = b/k(0) and b0 (0) = a.
ε1 (k(t2 ))2 (∂t21 + 1)v1 = −2k∂t1 ∂t2 v0 − k ′ ∂t1 v0 = −(a0 k ′ + 2ka′0 ) cos(t1 ) + (b0 k ′ + 2kb′0 ) sin(t1 ). (401)
To prevent secular terms, let a0 (t2 )k ′ (t2 ) + 2k(t2 )a′0 (t2 ) = b0 (t2 )k ′ (t2 ) + 2k(t2 )b′0 (t2 ) = 0. Thus,
da0 (t2 ) k ′ (t2 ) da0 1 k′ 1 dk 1 α0

=− a0 (t2 ) =⇒ =− dt2 = − =⇒ ln(a0 ) = − ln(k) + ln(α0 ) =⇒ a0 = √ . (402)
dt2 2k(t2 ) a0 2k 2 k 2 k
Follwing a similar procedure for b0 (t2 ) we arrive at
α0 β0
a0 = √ and b0 = √ , (403)
k k
where α0 and β0 are constants determined from the initial conditions. Our first-term approximation is therefore
Z t Z t
1
y∼ p α0 sin k(ετ )dτ + β0 cos k(ετ )dτ . (404)
k(εt) 0 0
p
The values of α0 and β0 are determined from the initial conditions, and one finds that β0 = y(0) k(0) and α0 =
−1/2 √
y ′ (0) [k(0)] . For the particular case a = 0, b = 1 we thus have β0 = 0, α0 = 1/ 2. p
As a check on the accuracy of the multiple-scale approximation in (404), the amplitude function A(t) = α0 / k(εt)
is shown in Fig.43. This is plotted, rather than the entire approximation, because of the highly oscillatory nature of the
solution. It is evident that the multiple-scale approximation gives a very accurate description of the solution.
There are other ways to derive the approximation in (404). One that is relatively easy to apply is the WKB method.
The advantage of the multiple-scales approach is that it works on nonlinear problems. This is evident in the original paper
on this problem by Kuzmak (1959), but this method requires us to be able to recognise the appropriate scales and then
determine what terms are responsible for secular terms.
6.7 Boundary Layers

As stated earlier, the method of multiple scales is not limited to the time variable. To illustrate this, we consider one of the
boundary-layer problems from the last chapter. Specifically, we take the following problem from 5.3:
εy ′′ + 2y ′ + 2y = 0, for 0 < x < 1, (405)

where
y(0) = 0 and y(1) = 1. (406)

We set x1 = x/ε and x2 = x, and equation (405) becomes

∂x21 + 2ε∂x1 ∂x2 + ε2 ∂x22 v + 2 (∂x1 + ε∂x2 ) v + 2εv = 0. (407)
2 Or d
R τ =x
in reverse, recalling that dx
f (τ )dτ = f (x) we see that:
τ =0
Z
τ =t Z Z Z
1 τ =t 1 ετ =εt 1 s=εt
f (t, ε) = k(ετ )dτ = k(ετ )d(ετ ) = k(ετ )d(ετ ) = k(s)ds
τ =0 ε τ =0 ε ετ =0 ε s=0
and so
df df dεt d2 f
= = k(εt) =⇒ 2 = εk ′ (εt).
dt d(εt) dt dt
92
It is worth noting the intervals that are used for the space scales we have introduced. For the boundary-layer scale, we have
that 0 ≤ x1 < ∞, and for the other coordinate we require 0 ≤ x2 ≤ 1. Now, assuming
v ∼ v0 (x1 , x2 ) + εv1 (x1 , x2 ) · · · , (408)

(with the requirement that this be asymptotic for x2 = Os (1)) one finds, from (407), that the general solution for the first
term is
v0 = a0 (x2 ) + b0 (x2 )e−2x1 . (409)

To determine the coefficients, we need the O(ε) equation from (407), and we find that

∂x21 + 2∂x1 v1 = −2∂x1 ∂x2 v0 − 2∂x2 v0 − 2v0 . (410)
The general solution of this is
v1 = a1 (x2 ) + b1 (x2 )e−2x1 − (a′0 + a0 )x1 + (b0 − b′0 )x1 e−2x1 . (411)
Since 0 ≤ x1 < ∞, to remove the secular terms we require a′0 + a0 = 0. In fact if we consider the last two O(ε) terms, we
see that there contribution to the solution will be
−2x2
(a′0 + a0 )εx1 + (b0 − b′0 )εx1 e−2x1 = (a′0 + a0 )x2 + (b0 − b′0 )x2 e ε (412)
and this term is no longer formally O(ε) As this would spoil the well-ordering of the asymptotic series, we require that
a′0 + a0 = 0 and b0 − b′0 = 0.3
Using the fact that v0 (0, 0) = 0 and v0 (∞, 1) = 1, it follows that the first-term approximation of the solution is
y ∼ e1−x − e1+x−2x/ε . (413)

In comparison, using matched asymptotics, we obtained using the composite expansion (188),
y ∼ e1−x − e1−2x/ε .
These expressions are not exactly the same but are asymptotically equivalent. The difference is in the second exponential
in (413). However, this term is exponentially small except in the boundary layer. Since the layer has thickness O(ε), it follows
that any contribution the x makes is of second order. In other words, the two expansions are asymptotically equivalent (to
first order). As a final comment, note that when using multiple scales, the balancing needed to determine the correct space
scale, and the location of the boundary layer, are hidden but are still a very important component of the solution procedure.
Those who may want to pursue the connections between matched asymptotic expansions and multiple scales should consult
the papers by Wollkind (1977) and Bouthier (1984).
6.8 Partial Differential Equations

To begin the study of the application of multiple scales to partial differential equations, we consider the equation for an
elastic string with weak damping. Specifically, we will study the problem
∂x2 u = ∂t2 u + ε∂t u, for 0 < x < 1 and 0 < t, (414)

where
u=0 for x = 0, 1, (415)

and
u(x, 0) = g(x) and ∂t u(x, 0) = 0. (416)

Even though a partial differential equation is involved, (414) is similar to the oscillator problem considered in section 6.5.
This is because the input is through the initial conditions (416), and there is weak damping. Consequently, it should not be
surprising to find out that a regular expansion results in secular terms.
Regular expansion Writing u ∼ u0 (x, t) + εu1 (x, t) we obtain the following problems:
ε0 : ∂x2 u0 = ∂t2 u0 ; u0 (x = 0, t) = u0 (x = 1, t) = 0; u0 (x, 0) = g(x),∂t u0 (x, 0) = 0.

This is a standard PDE problem which can be solved using Fourier series. At the next order we have:
ε0 : ∂x2 u1 = ∂t2 u1 + ∂t u0 ; u1 (x = 0, t) = u1 (x = 1, t) = 0; u1 (x, 0) = 0,∂t u0 (x, 0) = 0.

This is a linear inhomogeneous PDE. (Finish as exercise and demonstrate the existence of secular terms).
3 Alternatively, it is possible to show that in order to minimize the contribution of the second term, we must take b0 − b′0 = 0 (Exercise).
93
Multiple scales approach To use multiple scales, we need to determine the appropriate time scales. Given the similarity
with the weakly damped oscillator, it seems reasonable to assume that the same time scales can be used in this problem,
and so we take t1 = t and t2 = εt. Thus we transform
u(x, t) → v(x, t1 , t2 ) (417)

in the usual way using the chain rule, we find in this case that (414) becomes

∂x2 v = ∂t21 + 2ε∂t1 ∂t2 + ε2 ∂t22 v + ε (∂t1 + ε∂t2 ) v, (418)
where
v=0 for x = 0, 1, (419)

and
v = g(x) and (∂t1 + ε∂t2 ) v = 0 at t1 = t2 = 0. (420)

As before, the solution is not unique, and we will use this freedom to minimize the contribution of the second term in
the expansion. This will prevent the appearance of secular terms.
The appropriate expansion of the solution is
v ∼ v0 (x, t1 , t2 ) + εv1 (x, t1 , t2 ) + · · · . (421)

Introducing this into (418) yields the O(1) problem
∂x2 v0 = ∂t21 v0 , (422)

where
v0 = 0 for x = 0, 1, (423)
and
v0 (x, 0, 0) = g(x) and ∂t1 v0 (x, 0, 0) = 0. (424)

Noting that (422) is a pseudo PDE wrt to the variable t2 , using separation of variables, one finds that the general solution
of (422), (423) is
∞
X
v0 (x, t1 , t2 ) = [an (t2 ) sin(λn t1 ) + bn (t2 ) cos(λn t1 )] sin λn x, (425)
n=1
where λn = nπ. The initial conditions in (424) will be imposed once we determine the coefficients an (t2 ) and bn (t2 ).
The O(ε) equation that comes from introducing (421) into (418) is
∂x2 v1 = ∂t21 v1 + 2∂t1 ∂t2 v0 + ∂t1 v0

∞
X (426)
= ∂t21 v1 + An (t1 , t2 ) sin λn x,
n=1
where
An = −(2a′n + an )λn cos(λn t1 ) + (2b′n + bn )λn sin(λn t1 ).

Given the boundary conditions (419) we require v1 = 0 at x = 0, 1 and it is appropriate to expand the solution of (426)
in a modal expansion like the one that appears in (425). Thus, the solution of (426) has the form
∞
X
v1 = un (t1 , t2 ) sin λn x (427)
n=1
where un (t1 , t2 ) are the usual Fourier coefficients. Substituting this into (426) and using (425), one finds that
∂t21 un + λ2n un = −(2a′n + an )λn cos(λn t1 ) + (2b′n + bn )λn sin(λn t1 ).

Since the trigonometric functions on the right-hand side are solutions of the associated homogenous equation, to prevent
secular terms in the expansion we must take
2a′n + an = 0 and 2b′n + bn = 0.
94
Solving these equations, and using the initial conditions (424), one finds that an = 0 and bn = βn e−t2 /2 , where
Z 1
bn (0) = βn = 2 g(x) sin(λn x)dx
0
using Fourier series.
Therefore, our first-term approximation of the solution is
∞
X
u(x, t) ∼ βn e−εt/2 cos(λn t) sin λn x. (428)
n=1
It appears from this example that the application of the method of multiple scales to partial differential equations is
a relatively straightforward generalization of what was used for ordinary differential equations. It should be understood,
however, that this was a very simple example and partial differential equations can easily produce different, and very
interesting, uses of multiple scales.
6.8.1 A linear wave problem

To investigate how multiple scales can be used for travelling waves (f (x − ct)), we will again start with the equation for an
elastic string with weak damping. The problem is
∂x2 u = ∂t2 u + ε∂t u, for − ∞ < x < ∞ and 0 < t, (429)

where
u(x, 0) = F (x) and ∂t u(x, 0) = 0. (430)
Preamble: d’Alembert’s solution For the simple linear wave equation utt + c2 uxx = 0, we recall that if we introduce
the characteristic variables:
θ1 = x − ct; θ2 = x + ct; u(x, t) −→ v(θ1 , θ2 )
then the chain rule leads to:
∂u ∂v ∂θ1 ∂v ∂θ2
= + = vθ1 + vθ2
∂x ∂θ1 ∂x ∂θ2 ∂x
and
∂u ∂v ∂θ1 ∂v ∂θ2
= + = −cvθ1 + cvθ2
∂t ∂θ1 ∂t ∂θ2 ∂t
and the wave equation reduces to:
∂2v
4c2 = 0.
∂θ1 ∂θ2
This can be written as
∂ ∂v
= 0.
∂θ1 ∂θ2
∂v
As ∂θ 2
is at most a function of θ1 and θ2 , we can first regard this equation as being a first order pseudo PDE. Thus we
integrate wrt θ1 keeping θ2 constant. This results in
∂v
= a(θ2 )
∂θ2
where a(θ2 ) is an arbitrary function. We now regard this as a first order pseudo PDE and we integrate accordingly to get:
Z
∂v
= a(θ2 ) =⇒ v = a(θ2 )dθ2 + b(θ1 ) = c(θ2 ) + b(θ1 )
∂θ2
where c, b are arbitrary functions to be determined from the boundary conditions. In terms of the original variables this
solution takes the form:
u(x, t) = b(x − ct) + a(x + ct).
Note that this is a solution for any continuously differentiable functions a, b.
95
Regular expansion To see what happens when the multiple-scales method is not used, we try the regular expansion,
u ∼ u0 (x, t) + εu1 (x, t) + · · · . (431)

Substituting this into (429) yields the following:
ε0 ∂x2 u0 = ∂t2 u0 , u0 (x, 0) = F (x), ∂t1 u0 (x, 0) = 0. (432)

It is convenient at this point to switch to characteristic coordinates (d’Alembert’s technique) and let θ1 = x − t and
θ2 = x + t. In this case, ∂x = ∂θ1 + ∂θ2 , ∂t = −∂θ1 + ∂θ2 , and t = 0 corresponds to θ1 = θ2 . Also, the equation for u0 reduces
to ∂θ1 ∂θ2 u0 = 0. Integrating this yields the general solution
u0 = f0 (θ1 ) + g0 (θ2 ),
where f0 (θ1 ) and g0 (θ2 ) are determined from the initial conditions. From these conditions, one finds that the solution of the
O(1) problem is
1 1
u0 = [F (θ1 ) + F (θ2 )] = [F (x − t) + F (x + t)] .
2 2
ε1 (∂x2 − ∂t2 )u1 = ∂t u0 (433)

If we use characteristic coordinates to find the O(1) solution, the equation becomes
1 1
4∂θ1 ∂θ2 u1 = − F ′ (θ1 ) + F ′ (θ2 ).
2 2
Integrating this, we obtain
1 1
u1 = − θ2 F (θ1 ) + θ1 F (θ2 ) + f1 (θ1 ) + g1 (θ2 ),
8 8
where f1 (θ1 ) and g1 (θ2 ) are arbitrary functions. Since −∞ < x < ∞, this solution contains secular terms; from what we
have done so far, there is no way to prevent this from occuring.
Multiple scales approach Because of the appearance of secular terms in the regular expansion, we will use multiple
scales to find an asymptotic approximation of the solution. Unlike the problems examined earlier, this one appears to have
secular terms both in space and in time. Consequently, there are several scales we can introduce. For example, we can use
one of the following:
1. t1 = t, t2 = εt, x1 = x, x2 = εx,
2. θ1 = x − t, θ2 = x + t, φ1 = εθ1 , φ2 = εθ2 ,
3. θ1 = x − t, θ2 = x + t, t2 = εt, x2 = εx.
These are equivalent, but note that each list contains four variables. Unfortunately, this is usually necessary and means
that there is going to be some serious symbol pushing ahead.
We will use the scales given in option 3 above, which means that (using the chain rule)
∂x → ∂θ1 + ∂θ2 + ε∂x2 and ∂t → −∂θ1 + ∂θ2 + ε∂t2

and
u(x, t) → v(θ1 , θ2 , x2 , t2 ). (434)

Substituting these into (429), we obtain the following equation:

∂θ1 ∂θ2 + 2ε [∂θ1 + ∂θ2 ] ∂x2 − (−∂θ1 + ∂θ2 ) ∂t2 ] +O(ε2 ) v = ε [−∂θ1 + ∂θ2 + O(ε)] v. (435)
Also note that the initial conditions at t = 0 now correspond to θ1 = θ2 = x and t2 = 0. In fact, the initial conditions in
(430) become
v = F (θ1 ) and (−∂θ1 + ∂θ2 + ε∂t2 )v = 0 when θ1 = θ2 and t2 = 0. (436)

The next step is to introduce the expansion for the solution. For this problem, it is simply
v ∼ v0 (θ1 , θ2 , x2 , t2 ) + εv1 (θ1 , θ2 , x2 , t2 ) + · · · . (437)

This leads to the following problems:
96
ε0 4∂θ1 ∂θ2 v0 = 0, v0 (θ1 , θ1 , x2 , 0) = F (θ1 ), (−∂θ1 + ∂θ2 )v0 (θ1 , θ1 , x2 , 0) = 0. (438)
The general solution of this problem (ignoring the ICs for the moment) is
v0 = f0 (θ1 , t2 ) + g0 (θ2 , t2 ). (439)

These functions do not depend on the variable x2 . We are assuming that v0 does not depend on x2 because of the differential
equation and initial conditions. This hypothesis is made to reduce the effort to solve the problem, and the dependence can
always be restored later if necessary.
At the next order we find that:
ε1 4∂θ1 ∂θ2 v1 = 2(−∂θ1 + ∂θ2 )∂t2 v0 + (−∂θ1 + ∂θ2 )v0 .
Substituting the O(1) solution into this equation and solving, one finds that
1 1
v1 = − θ2 (2∂t2 + 1)f0 + θ1 (2∂t2 + 1)g0 + f1 (θ1 , t2 ) + g1 (θ2 , t2 ).
4 4
To prevent secular terms, we take
(2∂t2 + 1)f0 = 0 and (2∂t2 + 1)g0 = 0.

Integrating these equations and using the O(1) initial conditions yields
1 1
f0 =F (θ1 )e−t2 /2 and g0 = F (θ2 )e−t2 /2 . (440)
2 2
Therefore, a first-term approximation of the solution that is valid for large t is
1
u(x, t) ∼ [F (x − t) + F (x + t)] e−εt/2 . (441)
2
6.9 Other multiple scale examples

6.9.1 Forced Motion Near Resonance
The problem to be considered corresponds to a damped, nonlinear oscillator that is forced at a frequency near resonance. In
particular, the problem that is investigated is
y ′′ + ελy ′ + y + εky 3 = ε cos(1 + εω)t, for 0 < t, (442)

where
y(0) = 0 and y ′ (0) = 0. (443)

′ 3
In the differential equation, the damping (ελy ), nonlinearity (εky ), and forcing (ε cos(1 + εω)t) are small. Also, ω, λ, and k
are constants with λ and k non-negative. This oscillator has been studied extensively and is an example of a forced Duffing
equation (see Exercises). Note that if ε = 0 the equation becomes y ′′ + y = 0. For ε ≪ 1, the right hand side is thus small
and close to being resonant. Because of the small forcing, and the zero initial conditions, it is natural to expect the solution
to be small.
Simple example For example, if we were to consider the simpler equation
y ′′ + y = ε cos Ωt. (444)

then the exact solution is, assuming Ω 6= ±1,
ε
y(t) = [cos Ωt − cos t] . (445)
1 − Ω2
Therefore, the solution is small since it is O(ε). This observation is not true when the driving frequency Ω ≈ 1. To investigate
this situation, suppose Ω = 1 + εω. A particular exact solution in this case is, for ω 6= 0, −2/ε,
−1
y= cos(1 + εω)t; (446)
ω(2 + εω)
and for ω = 0 or ω = −2/ε, (Ω = ±1)
1
y= εt sin t. (447)
2
97
0.6
0.3
Solution
0
−0.3
−0.6
0 50 100 150 200 250
t - axis
Figure 44: Dashed curve shows the numerical solution of (442) with ε = 0.05, λ = 2, κ = 4, ω = 3/8. The solid curve is the
numerical solution of (461),(462) for the amplitude function A(εt) determined from the multiple scales approximation (464).
In both cases a relatively small, O(ε), forcing results in at least an O(1) solution. Also, the behaviour of the solution is
significally different for ω = 0 or ω = −2/ε than for other values of ω. This is typical of a system that is driven at one of its
characteristic frequencies, and it is associated with the phenomenon of resonance.
Such behaviour is also present in the nonlinear oscillator equation in (442). To demonstrate this, the values of y(t) obtained
from the numerical solution of (442) are shown in Fig.44. It is seen that the amplitude of the solution is an order of magnitude
larger than the forcing.
Another observation that can be made from Fig.44 is that there appear to be at least two time scales present in the
problem. One is associated with the rapid oscillation in the solution and the other is connected with the relatively slow
increase of the amplitude. Thus, the multiple-scales method is a natural choice for finding an asymptotic approximation to
the solution. We will try to give ourselves some flexibility and take t1 = t and t2 = εα t, where α > 0. In this case (442)
becomes

∂t21 + 2εα ∂t1 ∂t2 + ε2α ∂t22 v + ελ (∂t1 + εα ∂t2 ) v + v + εkv 3 = ε cos(t1 + εωt1 ). (448)
As stated earlier, because of the zero initial conditions (443) and the small forcing, one would expect the first term in
the expansion of the solution to be O(ε). However, near a resonant frequency the solution can be much larger. It is not clear
what amplitude the solution actually reaches, and for this reason we will take
v ∼ εβ v0 (t1 , t2 ) + εγ v1 (t1 , t2 ) + · · · , (449)

where β < 1 and β < γ. Substituting (449) into (448) yields the equation
(∂t21 + 2εα ∂t1 ∂t2 + ...)(εβ v0 + εγ v1 + ...) + ελ(∂t1 + ...)(εβ v0 + ...) (450)
1a 2a 1 (451)
+ εβ v0 + εγ v1 + ... + εκ(ε3β v03 + ...) = ε cos(t1 + εωt1 ) (452)
2a 2 3 (453)
expansion. With this, and (443), the following results:
εβ : (∂t21 + 1)v0 = 0, (454)

v0 (0, 0) = ∂t1 v0 (0, 0) = 0. (455)
The general solution of this problem is
98
v0 = A(t2 ) cos(t1 + θ(t2 )), (456)
where
A(0) = 0. (457)
The compact form of the solution is used in (456), rather than the form in (375), since it is easier to deal with the nonlinear
terms that appear in the next problem. In particular, we will make use of the identity cos3 φ = 14 [3 cos(φ) + cos(3φ)] .
We need to determine α, β, and γ before going any farther. To do this, note that the higher-order terms in expansion
(449), such as εγ v1 , are there to deal with the nonzero terms in (450). This includes the terms labeled 1 , 2 , and 3 . We
can take care of all three if they balance with the terms labeled with a 2a . This gives us γ = 1, 1 + β = γ, and 1 + 3β = γ.
The remaining observation to make is that all three terms 1 , 2 , and 3 will produce secular terms in the expansion, and
therefore we must choose our slow time scale t2 to deal with this. This means term 1a must balance with 1 , 2 , and 3 ,
which gives us α + β = 1. From these equations, we find that β = 0 and α = γ = 1. With this, the right-hand side of (450)
is rewritten as ε cos(t1 + ωt2 ) (because α = 1 and so t2 = εt). Thus, the next order problem to solve is the following:
ε1 : (∂t21 + 1)v1 = −2∂t1 ∂t2 v0 − λ∂t1 v0 − kv03 + cos(t1 + ωt2 ) (458)

k
= −(2A′ + λA) sin(t1 + θ) + 2θ′ A cos(t1 + θ) − A3 [3 cos(t1 + θ) + cos 3(t1 + θ)] + cos(t1 + ωt2 ). (459)
4
We want to determine which terms on the right-hand side of this equation produce secular terms - in other words, which
terms are solutions of the associated homogeneous equation (multiples of cos t1 and sin t1 ). This can be done by noting
cos(t1 + ωt2 ) = cos(t1 + θ − θ + ωt2 ) = cos(t1 + θ) cos(θ − ωt2 ) + sin(t1 + θ) sin(θ − ωt2 ). (460)
Note that if cos(t1 ) is resonant then so too is cos(t1 + θ) where θ is independent of t1 . This follows because cos(t1 + θ) =
cos θ cos t1 − sin θ sin t1 .
Thus, to remove the sin(t1 + θ) and cos(t1 + θ) terms in (458), it is required that
2A′ + λA = − sin(θ − ωt2 ) (461)

and
3k 3
2θ′ A − A = − cos(θ − ωt2 ). (462)
4
From (461), (462), the initial condition in (457), and assuming A′ (0) > 0, it follows that
π
θ(0) = − . (463)
2
To summarize our findings, the first-term expansion of the solution of (442) has the form
v ∼ A(t2 ) cos(t1 + θ(t2 )). (464)

To find the amplitude A(t2 ) and phase function θ(t2 ), it remains to solve (461), (462) along with initial conditions (457),
(463). Unfortunately, the solution of this nonlinear system of equations is not apparent. In lieu of the exact solution, we can
investigate the general properties of this sort of system and, in particular, steady state solutions.
Steady state solution

In particular, we will investigate what value the amplitude A approaches if it goes to a steady state as t → ∞ (if this limit
exists). So, assuming A′ = 0, let A = A∞ , where A∞ is a positive constant. From (461), it follows that θ does not reach a
steady state in this situation because if θ = θ∞ when A = A∞ , A′∞ = 0 , then (461) tells us that
λA∞ = − sin(θ∞ − ωt2 ) (465)

and this equation makes no sense because t2 is a variable. For the equation to make sense we require that θ − ωt2 = θ0 for
some constant θ0 when A = A∞ where θ0 is a constant that satisfies sin(θ0 ) = −λA∞ . Thus when A = A∞ and θ = ωt2 + θ0
we have θ′ (t2 ) = ω. Thus (461) and (462) become
λA∞ = − sin(θ − ωt2 ) (466)

3k 2
2ωA∞ − A = − cos(θ − ωt2 ). (467)
4 ∞
99
Aoo
AM
||
| Am |||
ωm ωM ω
Figure 45: A graph of the A∞ as a function of ω showing the positive solutions of (468) when λ, κ > 0 (the exact shape of
the curve depends on λ, κ). The curve has three branches: I for −∞ < ω ≤ ωm , II for ωm < ω < ωM , III for ωM ≤ ω < ∞.
Note that AM = 1/λ and ωM = 3κ/(8λ2 ).
As cos2 + sin2 = 1 we see that the constant A∞ is determined from the algebraic equation
" 2 #
3k
A2∞ λ2 + 2ω − A2 = 1. (468)
4 ∞
In principle λ, k, ω are fixed constants. But now it is interesting to see the effect of different values of ω, which corresponds
to varying the frequency of the forcing term. An example of the possible (obtained numerically) steady state solutions (A∞ )
obtained from (468) is shown in Fig.45 as a function of the frequency parameter ω. The interesting parts of the curve occur
for ωm < ω < ωM , where there are three possible steady states for A. The middle one, branch II, is unstable and cannot
be reached unless the initial conditions are chosen so that the solution starts exactly at one of these points (the justification
of this statement is beyond the scope of this course). The other two branches I and III, are stable, and the solution will
approach one of them depending on initial conditions. This situation is known as bistability. Note that if one were to start on
branch I and then increase the driving frequency to the point that ω > ωM , then the amplitude would jump down to branch
III. Similarly, once it is on the lower part of this solution curve and the frequency is decreased below ωm , the amplitude then
jumps up to branch I.
A comparison between the amplitude function A(t2 ) obtained via numerical solution of (461), (462) and the full numerical
solution of (442) is shown in Fig.44. Note that though (461), (462) can only be solved numerically in general, it is much
easier to do so than to numerically solve the full problem (442). This type of comparison is made, rather than plotting all
the values for the solution that are obtained from (464), because of the oscillatory nature of the solution. In any case, it is
clearly seen that the agreement is quite good. However, one might argue that if the reduced problem (442) must be solved
numerically then why not just solve the original problem numerically? It is certainly true that the numerical solution is not
limited to small values of the parameter ε. In this regard it is superior to the asymptotic approximation. The asymptotic
approximation, on the other hand, is capable of providing quick insight into how the solution depends on the various
parameters of the problem. This sort of pro-and-con argument applies to almost all numerical/asymptotic solutions. What
is special about the present example is the oscillatory nature of the response. In this case, the asymptotic approximation has
the rapid oscillations built into the representation. Thus, the functions that need to be computed numerically are relatively
smooth. In this sense, it is easier to solve (461), (462) numerically than the original equation (442).
The analysis here has been for what is sometimes called harmonic, or primary, resonance. There are also subharmonic and
superharmonic resonances. To see where these originate, note that the solution in (445) is periodic if Ω = n or if Ω = 1/n,
where n is a positive integer with n 6= 1. For the first case, the solution has period 2π while the forcing has period 2π/n. This
corresponds to what is known as subharmonic oscillation. If Ω = 1/n, then the solution and forcing have period 2εn. This
is known as superharmonic, or ultraharmonic, oscillation. With the nonlinear Duffing equation (442), it is possible for there
to be a resonant response if Ω is equal to subharmonic or superharmonic frequency. However, the amplitudes are generally
smaller than what is found for harmonic resonance. Multiple scales can be used to analyze the subharmonic response.
6.9.2 Nonlinear Waves

Traveling wave solutions are found for a wide variety of nonlinear partial differential equations. This differs significantly from
the situation for linear problems where wave propagation is limited to particular types of equations. In this section, two
examples are analyzed that illustrate some of the situations where nonlinear equations are possible and how to use multiple
scales to find them.
100
Example: Nonlinear Wave Equation Consider the nonlinear wave equation
∂x2 u = ∂t2 u + K(x, u, ε), for − ∞ < x < ∞ and 0 < t, (469)
where
u(x, 0) = F (x) and ∂t u(x, 0) = G(x). (470)
Preamble There is little hope of solving this problem in closed form, so our approach is to investigate particular
solutions that are in the form of traveling waves. In the example to follow, we will examine the behaviour of the solution of
the non-linear Klein-Gordon equation, where K(x, u, ε) = u + εu3 and the equation becomes:
∂x2 u = ∂t2 u + u + εu3 . (471)

This type of equation comes up in a variety of situations. For example, it describes the motion of a string on an elastic
foundation as well as the waves in a cold electron plasma, and it also arises in quantum mechanics.
To explain the approach we will take, consider the linear equation that is obtained when ε = 0. In this case, rather than
use d’Alembert’s approach, as the equation is linear with constant coefficients, we can expect exponential type solutions and
we try the (plane wave, periodic wave train) substitution
u = exp(i(kx − ωt)) (472)

and we find that
ux = iku; uxx = −k 2 u; ut = −iωu; utt = −ω 2 u
so that (471) with ε = 0 reduces to:
−k 2 u = −ω 2 u + u =⇒ ω 2 = k 2 + 1
i.e., (472) is a solution of the equation for any value of k provided ω 2 = 1 + k 2 . As we are working on an infinite domain,
one could use the Fourier transform and find a general solution of the linear equation uxx = utt + u in the form
Z ∞ Z ∞
u(x, t) = A(k)ei(kx−ωt) dk + B(k)ei(kx+ωt) dk,
−∞ −∞
where A(k) and B(k) are determined from the initial conditions and ω = (1 + k 2 )1/2 . It is apparent from this that the
solution is composed of the superposition of the plane wave solutions
u± (x, t) = ei(kx±ωt) = cos(kx ± ωt) + i sin(kx ± ωt).

It is therefore worth investigating what effect the nonlinear term εu3 has on one of these wave planes. We will consider
u(x, t) = cos(kx − ωt), using an expansion where this function is effectively the first term, and we will then determine how
the nonlinearity contributes to the second term.
Boundary conditions Thus, motivated by the above discussion we will consider (471) subject to the initial conditions
derived from the exponential solution exp(i(kx − ωt))
u(x, 0) = α cos kx; ut (x, 0) = αω sin kx (473)

where α is an given initial amplitude.
To find an expansion of the solution, the simplest thing to try is one of the form
u(x, t) ∼ u0 (kx − ωt) + εu1 (x, t) + · · · . (474)

However, it is not hard to show that this leads to secular terms. Consequently, we will use multiple scales, and this brings up
choices for the scales similar to those we had for the linear wave problem see (6.8.1). Our first guess might be to introduce
scales x1 = x, t1 = t, x2 = εx, and t2 = εt. Because we are investigating what effects the nonlinearity has on the right running
wave, for the scales we can simplify matters by taking
θ = kx − ωt, x2 = εx, t2 = εt, u(x, t) → v(θ, x2 , t2 ).

In this development, we assume that ω = (1 + k 2 )1/2 and k > 0. The differential equation in this case becomes

∂θ2 − 2ε(k∂x2 + ω∂t2 )∂θ + O(ε2 ) v + v + εv 3 = 0. (475)
The appropriate expansion of the solution is now
v(x, t) ∼ v0 (θ, x2 , t2 ) + εv1 (θ, x2 , t2 ) + · · · . (476)
101
This leads to the following problems
O(1) (∂θ2 + 1)v0 = 0.

The general solution of this
v0 = A(x2 , t2 ) cos [θ + φ(x2 , t2 )] . (477)

As usual, the functions A and φ are used to prevent secular terms from appearing in the higher-order terms in the expansion.
The initial conditions will be discussed after the general form of the expansion has been determined. At next order we find
that
1
ε1 (∂θ2 + 1)v1 = 2(k∂x2 + ω∂t2 )∂θ v0 − v03 = −2 [(k∂x2 + ω∂t2 )A] sin(θ + φ) − A3 cos 3(θ + φ) (478)
4
3 2
− 2 (k∂x2 + ω∂t2 )φ + A A cos(θ + φ) (479)
8
This is a pseudo PDE much like what we have seen earlier. To prevent secular terms, it is required that
(k∂x2 + ω∂t2 )A = 0 (480)

and
3
(k∂x2 + ω∂t2 )φ + A2 = 0. (481)
8
These first-order equations can be solved using characteristic coordinates. Thus let r = ωx2 + kt2 and s = ωx2 − kt2 and we
now consider A = A(r, s), φ = φ(r, s). With this (480), (481) become
3
∂r A = 0 and ∂r φ = − A2
16ωk
3
and we see that A must be independent of r i.e., A = A(s) and φ = − 16ωk A2 r + φ0 (s) with the actual functions A, φ0 still
to be determined. Thus a first-term approximation of the solution of the nonlinear Klein-Gordon equation is

3
v ∼ A cos θ − (ωx2 + kt2 )A2 + φ0 , (482)
16ωk
where A(x2 , t2 ) ≡ A(ωx2 − kt2 ) and φ0 (x2 , t2 ) ≡ φ0 (ωx2 − kt2 ) are arbitrary.
We began this example by asking what the nonlinearity does to the plane wave solution. To answer this, we impose the
initial conditions (473) corresponding to the plane wave. In terms of the new multiple scales system used to derive (482),
these translate into the conditions
v0 (θ, x2 , 0) = α cos θ; ∂θ v0 (θ, x2 , 0) = −α sin θ.

3
In (482) this means A(ωx2 ) = α and φ0 (ωx2 ) = 16k A2 x2 . Thus, a first-term approximation of the solution in this case
(written in terms of the original variables u(x, t) ) is

3εα2
u(x, t) ∼ α cos kx − ω + t . (483)
16ω
Here α is the amplitude and so that the perturbed phase velocity

3εα2 ω
vph = 1 +
16ω 2 k
is increased by the nonlinearity and this increase is amplitude dependent. Specifically, the larger the amplitude, the faster
the wave moves. This dependence of the velocity on amplitude is typical of nonlinear waves and does not occur for linear
problems.
As it turns out, this example was one of the first wave problems to be analyzed using multiple scales (Luke, 1966). It has
since been applied to numerous other problems, such as nonlinear diffusion problems, which brings us to the next example.
102
Example: Diffusion equation with nonlinearity (including Fisher’s equation)
Population dynamics In the study of elementary population dynamics it is often proposed that a population is governed
by the logistics law, which states that the rate of change of population u = u(t) is given by:
du u
= ru 1 − (484)
dt K
where r > 0 is the growth rate and K > 0 is the carrying capacity. If u(t) is initially small, the linear growth term ru
might be expected to dominate the quadratic u2 term and the population would grow exponentially. As u becomes large, the
quadratic competition term −ru2 /K kicks in to inhibit the initial growth. Thus this corresponds to growth of a population
u when there is a limit on the size of the population K which the habitat can support; if u > K, then the RHS is negative
and the population decreases in size whenever u is greater than the limiting size. Note that for large times, one might expect
the population to equilibrate towards the asymptotically stable steady state solution u = K, the carrying capacity.
Now suppose that u is a population density i.e., population per unit volume and depends on the spatial variable x as well
as the time so that u = u(x, t). The conservation law may be formulated in the form:
ut + qx = ru(1 − u/K) (485)
where f = f (u) = ru(1 − u/K) is the assumed local source term given by the logistics growth law, and q is the population
flux or rate of flow. If we assume Fick’s law for the rate of flow i.e., q = −Dux then we arrive at the so-called Fisher’s
equation:
ut − Duxx = ru(1 − u/K). (486)

The RHS term is a source (or reaction) term and this is usually termed a reaction-diffusion equation.
Exercise Considering the dimensions of the parameters D, r, K, non-dimensionalise the Fisher equation to obtain the
dimensionless form
ut + uxx = u(1 − u). (487)

Note that two obvious solutions of this equation are u = 0 and u = 1. Both of these solutions are independent of t (i.e.,
steady state solutions) and independent of x (sometimes called uniform or spatially homogeneous solutions). We will also
see that this equation also (surprisingly) has travelling wave solutions.
Wave propogation In the study of wave propagation in chemical and biological systems, one comes across the problem
of having to solve
ut = εuxx + f (u), for − ∞ < x < ∞ and 0 < t, (488)

where
u(x, 0) = g(x). (489)

Note that if f (u) = u(1 − u) this can be considered a special case of (487) and in, fact, corresponds to the case of weak
diffusion when ε ≪ 1. This form can arise, for example, when there is also a given length scale L in the problem (originating
perhaps in the initial condition).
There are restrictions on the functions f (u) and g(x), and these will be given later. What is interesting is that the solution
of this nonlinear diffusion problem can evolve into a traveling wave of the form u(x, t) = U (x − vt), where v is the velocity of
the wave. Such a solution is impossible with linear diffusion, and this is one of the reasons why nonlinear diffusion problems
have received a lot of attention over the last few years. We will use multiple scales to find an asymptotic approximation of
the solution of (488), and our approach will be based on the work of Berman (1978). An introduction to the theory and
some of the applications of this equation can be found in Britton (1986) and Murray (1989).
We first need to list the restrictions made on the nonlinearity in (488). In doing so, we specify only what happens for
0 ≤ u ≤ 1. This interval for u(x, t) is the one of physical interest since this function is a nondimensional concentration
or density. In what follows we will assume that f (u) is smooth for 0 ≤ u ≤ 1 and positive for 0 < u < 1. Moreover,
f (0) = f (1) = 0, f ′ (0) > 0, and f ′ (1) < 0. Perhaps the simplest example, and the most well studied, is the quadratic
function f (u) = u(1 − u). This gives rise to what is known as Fisher’s equation or the Kolmogorov-Petrovsky-Piskunov
(KPP) equation. This has been used to model such phenomena as the spatial spread of certain genetic characteristics, the
propagation of cosmic rays in the atmosphere, and the evolution of a neutron population in a reactor. It is also worth
pointing out that this quadratic function is the same as appears in the logistic equation, which is studied in connection with
population dynamics.
103
The initial condition in (489) must also be consistent with the physics. The requirements on this function are that
0 ≤ g(x) ≤ 1, g(x) → 1 as x → −∞, and g(x) → 0 as x → ∞. On foot of the discussion above, in what follows we will simply
take
1
g(x) = , f (u) = u(1 − u) (490)
1 + eλx
where λ > 0 is a given constant.
The origin of all these conditions may seem mysterious, and it may be unclear what connections, if any, there are between
the choices we are making for f (u) and g(x). To explain what we have done, note that our requirements on f (u) result in
two steady-state solutions of the differential equation, namely u = 0 and u = 1. The function f (u) has also been set up so
that u = 1 is stable and u = 0 is unstable. Therefore, by requiring that g → 1 as x → −∞, and g → 0 as x → ∞, the initial
condition effectively connects these two steady states. The solution of the problem in this case is fairly simple since at any
given point on the x−axis, u(x, t) simply moves away from u = 0 and approaches u = 1. What is important is that the speed
at which it approaches u = 1 depends nonlinearly on the solution. The resulting movement toward u = 1 is what gives rise
to the traveling wave. This is a very simplistic explanation of what occurs, but hopefully it will provide some insight into
the behaviour of the solution.
Our assumption that ε is small in (488) means that there is weak diffusion, and this introduces a second time scale into
the problem. To account for this, we introduce the two time scales t1 = t and t2 = εt. The equation in (488) then takes the
form
(∂t1 + ε∂t2 )u = εuxx + f (u). (491)

There is nothing unusual about this equation, and so the appropriate expansion for the solution is a regular power series of
the form
u ∼ u0 (x, t1 , t2 ) + εu1 (x, t1 , t2 ). (492)

The O(1) problem that comes from substituting (492) into (491) and (489) is
∂t1 u0 = f (u0 ), (493)

where
u0 (x, 0, 0) = g(x). (494)

By separating variables in (493), one finds that
Z u0
dr
= t1 + θ(x, t2 ), (495)
1/2 (r)
f
where θ(x, t2 ) is arbitrary. Imposing the initial condition (494) yields
Z g(x)
dr
θ(x, 0) = . (496)
1/2 f (r)
Without knowing exactly what the function f (u) is, we will have to leave u0 (x, t1 , t2 ) defined implicitly in (495). However,
note that we can write the solution in the form u0 = U0 (s), where s = t1 +θ(x, t2 ). This form will prove handy when removing
secular terms.
To determine the second term in the expansion, note that the O(ε) equation obtained from (491) is

∂t1 u1 = f ′ (u0 )u1 + ∂x2 u0 − ∂t2 u0 = f ′ (u0 )u1 + f (u0 ) θxx − θt2 + f ′ (u0 )θx2 (497)
Equations (495) and (493) have been used in the last step. Fortunately, this equation for u1 is linear and is not hard to solve.
To do this, it is worth pointing out that, from (493), f ′ (u0 ) = ∂t1 ln(f (u0 )). Using this result, one finds that the solution of
the above equation for u1 is

u1 = A(x, t2 ) + t1 (θxx − θt2 ) + θx2 ln(f (u0 )) f (u0 ), (498)
where A(x, t2 ) is arbitrary. From a quick inspection of this solution, it would appear that secular terms are going to arise
unless we restrict θ(x, t2 ). Note, however, that u1 = o(u0 ) as t1 → ∞ or as θ → ∞. An example which illustrates this situation
very nicely is the case of when f (u) = u(1 − u) and the details are left to the reader. Problems associated with secularity
occur when θ → −∞. In this limit one finds that the functions u0 and f (u0 ) are comparable and ln(f (u0 )) ∼ k(t1 + θ) where
k = f ′ (0) is a positive constant. Thus, ln(f (u0 ) produces a secular term in (498) that is of the same order as the other t1
terms in that expression. Based on these observations, we require
104
θxx − θt2 + kθx2 = 0. (499)
This equation can be solved if we first make the change of variables w(x, t2 ) = ekθ(x,t2 ) . In this case (499) reduces to the
linear diffusion equation
wxx = wt2 , for − ∞ < x < ∞ and 0 < t2 . (500)

The solution of this is
Z ∞
1 √ 2
w(x, t2 ) = h(x + 2r t2 )e−r dr, (501)
π 1/2 −∞
where h(x) = w(x, 0). From (496) we get

Z !
g(x)
dr
h(x) = exp k . (502)
1/2 f (r)
In stating that the solution of (500) is given by the formula in (501), we are assuming that the function h(x) is well enough
behaved that the integral in (501) is defined for all x and t2 . The exponential term in (501) helps a lot, and because of it
2
we can get away with simply requiring that h(x) ≤ M eαx for −∞ < x < ∞, where α, M are constants with α < 1. This
condition usually is not particularly restrictive. However, in looking at how h(x) is defined in (502), it appears that we need
to be careful because f (u) is zero when u = 0 or u = 1 and g(x) approaches these values as x → ±∞. However, the conditions
we have imposed on f (u) and our assumed form of g(x) in (490) make it so this is not an issue. This will be demonstrated
later in the case of the quadratic nonlinearity found in Fisher’s equation. An examination of other possibilities can be found
in Larson (1978).
From the preceding analysis, we have found that a first-term approximation of the solution of (488) is u ∼ u0 (x, t1 , t2 )
where u0 is defined implicitly through the equation
Z u0
dr 1
= t1 + ln [w(x, t2 )] , (503)
1/2 f (r) k
where x(x, t2 ) is given in (501) and (502).
As stated earlier, nonlinear diffusion problems can give rise to travelling wave solutions. The asymptotic approximation
of the solution given in (503) can be used to demonstrate this.
Fisher’s equation To show how, we will take f (u) = u(1 − u) so that we are dealing with Fisher’s equation. In this case,
we have from (495) that u0 = U0 (t1 + θ(x, t2 )), where
1
U0 (s) = .
1 + e−s
This can be rewritten as
w(x, t2 )
u0 (x, t1 , t2 ) = , (504)
w(x, t2 ) + e−t1
where, from (501),
q Z ∞
t2
w(x, t2 ) = e−λx π e−t2 s(s−2λ) ds
−∞
−λx+λ2 t2
=e .
Therefore, u0 = U0 (t1 − λx + λ2 t2 ), and this gives us a traveling wave with velocity
1 + λ2 ε
v∼ . (505)
λ
What is interesting is that (505) shows that the speed of the wave depends on the shape of the initial profile. The steeper it
is (so λ increases), the slower it moves.
To illustrate the character of the traveling wave that comes from the nonlinear diffusion equation, the asymptotic solution
of Fisher’s equation is shown in Fig. 46. It is seen that the wave front moves to the right with a velocity of about 1/λ. Also
shown in Fig. 46 is the wave at t = 5 and at t = 10 as obtained from the numerical solution of the full problem. These
curves are included to demonstrate the accuracy of the multiple-scales approximation.
105
1.0
0.8
0.6
Solution
0.4
0.2
0
-10 -5 0 5 10 15 20
x - axis
Figure 46: Travelling wave solution of Fisher’s equation determined from the asymptotic solution in (504). The solution is
shown at t = 0, 5, 10. The numerical solution of the original problem (488) is shown with the dashed curves. ε = 0.01, λ = 1.
7 Method of strained co-ordinates

This method is due to Lighthill in 1949. It therefore predates the method of matched asymptotic expansions, and it is more
limited. Lighthill’s rule of thumb for this technique is that higher approximations should be no more singular than the
leading order approximation.
The method assumes that the function to be calculated f (x, ε) (i.e., the solution of the problem being considered) has
the same form as the unperturbed function f (x, 0), but the x-scale is slightly shifted and/or slightly distorted, i.e., a strained
co-ordinate. A small strain can make an asymptotic expansion for f (x, ε) non-uniform in x
• An example with a small distortion is
1 1 ε ε
f (x, ε) ≡ ∼ − 2 + 3 as ε → 0. (506)
x+ε x x x
This is not uniformly asymptotic: there is trouble at x = Os (ε) as ε → 0. For example for ε small but fixed ε/x2 ≫ 1/x
when x ≪ ε.
• Another example is
1
f (x, ε) ≡ sin (1 + ε)x ∼ sin x + εxcos x − ε2 x2 sin x as ε → 0.
2
This is not uniformly asymptotic: there is trouble at x = Os (ε−1 ) as ε → 0.
In both these examples, the leading order term is essentially the correct form: it just needs a slight shift or distortion to
make it uniformly asymptotic.
In solving differential equations we shall seek a near identity transformation x(s, ε) which makes f (x, 0) uniformly asymp-
totic to f (x, ε).
Thus we pose
f (x, ε) ∼ f0 (s) + εf1 (s),

x(s, ε) ∼ s + εx1 (s)
with the constraint that the expansion for f (x, ε) (but usually not for x(s, ε)) is uniformly asymptotic (i.e., we may allow
s ∼ εx1 (s) if necessary). This latter point can be important as if the series for x(s, ε) is not asymptotic (i.e., not necessarily
well ordered), it appears to break one of our fundamental principles. The point is that the freedom gained in allowing this
is used to ensure that the series for f (x, ε) is certainly asymptotic and well-ordered.
The excessive freedom in looking for the two functions in place of one is controlled by keeping everything simple in some
appropriate sense (quite subjective). Here we are essentially caught by the non-uniqueness of non-Poincaré expansions.
106
Before applying the method of strained co-ordinates one needs some reason to believe that f (x, ε) is like f (x, 0): the
method will not work when this is not true, such as in boundary layer problems.
The general rule in choosing our straining is that higher order approximations should be no more singular than the leading
order approximation (Lighthill). To clarify what we mean by “more singular”, note that in (506), for small enough x, the
second terms become larger (i.e., more singular) than the first. (Here we think of ε being as small as we like but fixed and
we then consider x → 0).
7.1 A model problem

In the model problem of Carrier we seek f (x, ε) as ε → 0 where f (x, ε) satisfies
(x + εf )f p + f = 1 in 0 ≤ x ≤ 1
subject to f = 2 at x = 1.
We can expect trouble near x = 0 where x fails to be a good approximation to (x + εf ).
We start by treating the problem naively as a regular perturbation problem. When the method of strained co-ordinates
is used, this naive expansion is always useful. A quick calculation yields.
1+x 3x2 − 2x − 1
f (x, ε) ∼ +ε .
x 2x3
The expansion is asymptotic at fixed x as ε → 0, but fails to be asymptotic at x = Os (ε1/2 ).
The exact solution can be obtained by integrating once and solving a quadratic for f :
1/2
x x2 2(x + 1)
f (x, ε) = − + + +4 .
ε ε2 ε
Expanding for x fixed as ε → 0 gives the above naive expansion.
Comparing the graph of the exact solution with that of our naive expansion, we see in figure 6.1 that the shapes are
similar, but that the exact solution has a ‘singularity’ at x = −εf whereas all the approximations have singularities at x = 0.
All we need to do is to ‘blow over’ the singularities of the approximations. This discussion suggests that the problem is one
suitable for the method of strained co-ordinates.
Figure 47: The continuous curve gives the exact solution for ε = 0.1. The short dashed curves on the right give the naive
approximation at leading order (upper curve) and with the ε−correction term (lower curve). The dash-dot curve is the
singular line x = −εf for the exact solution (the corresponding singular line for the naive approximation is the axis x = 0).
The dotted curve is the uniformly asymptotic approximation corresponding to the simplest choice for x1 in the method of
strained co-ordinates.
107
7.1.1 Solution by strained co-ordinates
We pose two expansions
f (x, ε) ∼ f0 (x) + εf1 (s)

x(s, ε) ∼ s + εx1 (s)
with f0 to be uniformly asymptotic over 0 ≤ x ≤ 1, which is roughly 0 ≤ s ≤ 1 but not precisely so.
In order to substitute into the governing equation, we need the change of variable in the differentiation,
d ds d ′ d
= = (1 − εx1 + · · · )
dx dx ds ds
Now substituting into the equation and the boundary condition and comparing coefficients of εn in the usual way, we find
′
ε0 : sf0 + f0 = 1 with f0 = 2 at s = 1.
Thus we recover the leading order result of the naive approximation,
1+s
f0 = .
s
At the next order we find that:
′
′ ′ x1 x 1 ′ 1
ε : (sf1 ) ≡ sf1 + f1 = −f0 (x1 + f0 − sx1 ) = 2 − 1 + 3 + 2 .
1 ′
(507)
s s s s
An unusual feature of this model problem is that we can write down the general solution to the above equation for f1 . In a
normal problem one does not need the general solution. Without actually solving (507) we see that one way to ensure that
f1 (s) is no more singular than f0 (s) is to demand that the RHS be zero. This would lead to:
′ c1
sf1 + f1 = 0 =⇒ f1 =
s
where c1 is set by the boundary conditions. Note that f1 is no more singular than f0 (Lighthill’s rule). We now see that we
can get the same result by examination of f1 .
A x1 1 1
f1 =
− 2 − 3 − 2.
s s 2s s
where A is the constant of integration. Now the boundary condition is to be applied at x = 1, i.e., at s = 1 − εx1 + O(ε2 ).
This gives the condition
′
f0 (1) + ε[−x1 (1)f0 (1) + f1 (1)] + O(ε2 ) = 2.
Therefore f1 (1) = −x1 (1), and so making our general solution satisfy the boundary condition we have

x1 3s2 − 2s − 1 1 2s + 1 3
f1 = − 2 + = − x1 + + .
s 2s3 s2 2s 2s
We now must select x1 so that f0 (s) is uniformly asymptotic to f (x, ε) on 0 ≤ x ≤ 1. Clearly we should remove the s−3 and
s−2 terms from f1 which are more singular that f0 = (1 + s)/s ∼ 1/s for small s. Thus the simplest choice is
1 + 2s 3
x1 = − with f1 (s) = .
2s 2s
With this choice for x1 , one can invert x = s + εx1 (s) for an approximation to s(x, ε) which can then be substituted into f0 .
This yields the uniformly asymptotic expansion
1/2
x x2 2(x + 1)
f (x, ε) ∼ − + 2
+ +1 .
ε ε ε
Comparing this expression with the exact solution, we see that the above approximation has a relative error Os (ε) uniformly
on 0 ≤ x ≤ 1. Adding in the εf1 (s) reduces the relative error to Os (ε2 ) for x fixed as ε → 0, but leaves the relative error as
Os (ε) for x =)s (ε1/2 ). To reduce the error when x = Os (ε1/2 ), we need to include ε2 x2 (s).
An alternative choice for x1 (s) is one that vanishes at the boundary, and consequentially leads to no forcing of f1 (s) from
the boundary condition.
3s2 − 2s − 1
x1 (s) = with f1 (x) ≡ 0.
2s
Due to the simplicity of the model problem, this choice gives the exact solution
108
Because we only have to make f0 (s) uniformly asymptotic to f (x, ε) in 0 ≤ x ≤ 1, which is not strictly 0 ≤ s ≤ 1 but
rather Os (ε1/2 ) ≤ s ≤ 1 + Os (ε), we can allow f1 (s) to be more singular than f0 (s). Thus a third choice for x1 just takes
out the worst singularity in f1 :
1 3s − 2
x1 (s) = − with f1 (2) = .
2s 2s2
This choice gives a fractional error Os (ε1/2 ) when x = Os (ε1/2 ), worse because of the dangerous choice of f1 .
Note there is no need to calculate f1 (s) if only the leading order uniformly asymptotic expansion, f0 (s) with x = s+εx1 (s),
is required.
One can choose an x1 which removes the forcing of that part of f1 which could be more singular than f0 (the first choice
above), or one can choose the x1 which makes f1 vanish identically (the second choice above), although the calculation of
this x1 involves more work.
7.1.2 Solution by method of matched asymptotic expansions

The naive expansion (f ∼ (1 + x)/x) serves as an outer approximation. It breaks down at x = Os (ε1/2 ) where f = Os (ε−1/2 ).
This suggests a rescaling
x = ε1/2 ξ
for the inner approximation and an expansion in powers of ε1/2 starting at ε−1/2 , i.e.,
f (x, ε) ∼ ε−1/2 F0 (ξ) + F1 (ξ) + ε1/2 F2 (ξ).
Substituting into the governing equation and comparing coefficients of εn/2 in the usual way we find the following.
′
ε−1/2 : (ξ + F0 )F0 + F0 = 0.
The general positive solution is
F0 = (ξ 2 + A0 )1/2 − ξ
with A0 a constant of integration.
′ ′ ′
ε0 : ξF1 + F1 F0 + F0 F1 + F1 = 1.
This linear equation has a general solution
ξ + A1
F1 =
(ξ 2+ A0 )1/2
with A1 a constant of integration.
Matching We will illustrate van Dyke matching at leading order. The outer solution is f ∼ (1 + x)/x to O(1) and the
inner solution is
1 1
f ∼ √ F0 = √ ((ξ 2 + A0 )1/2 − ξ)
ε ε
√ √
to O( ε). Note that in the matching we compare f and F/ ε not f and F .
The inner expansion of the outer solution is obtained by switching independent variable and expanding for small ε as far
as O(ε) with ξ = O(1). Thus
1+x 1 1 1
= 1 + √ ∼ √ (≡ ).
x εξ εξ x
The outer expansion of the inner solution is obtained by switching dependent variable and expanding for small ε as far as
O(1) with x = O(1). Thus r
1 2 1/2 x ε A0
√ ((ξ + A0 ) − ξ) = 1 + 2A0 2 ∼ .
ε ε x x
Thus clearly matching requires
A0 = 1
and we can form a composite solution in the usual way:
r
x 2 2 x
f∼ + − + 1.
ε ε ε
This completes the solution by matched asymptotic expansions at leading order. Further terms can be obtained in the usual
way.
Comparing the two methods of solution, we see that the method of strained co-ordinates produces a uniformly asymptotic
approximation in one step, as it were.
109
8 The WKB Method
8.1 Introduction
In the method of matched asymptotic expansions studied in chapter 5, the dependence of the solution on the boundary-
layer coordinate was determined by solving the boundary-layer problem. In a similar way, when using multiple scales, the
dependence on the fast time scale was found by solving a differential equation. This does not happen with the WKB method
because one begins with the assumption that the dependence is exponential. This is a reasonable expectation since many
of the problems we studied in 5 ended up having an exponential dependence on the boundary-layer coordinate. Also, by
making this assumption, we can significantly reduce the work necessary to find an asymptotic approximation of the solution.
The popularity of the WKB method can be traced back to the 1920s and the development of quantum mechanics. In
particular, it was used to find approximate solutions of Schrödinger’s equation. The name of the method is derived from
three individuals who were part of this development, namely, Wentzel, Kramers, and Brillouin. However, as often happens
in mathematics, the method was actually derived much earlier. Some refer to it as the method of Liouville and Green, who
both published papers on the procedure in 1837. It appears that this too is historically incorrect, since in 1817 Carlini used a
version of the approximation in studying elliptical orbits of planets. Given the multiple parenthood of the method, it should
not be unexpected that there are other names it goes by, including the phase integral method, the WKBJ method (the J
standing for Jefferys), the geometrical acoustics approximation, and the geometrical optics approximation. The history of
the method is surveyed very nicely by Heading (1962) and in somewhat more mathematical detail by Schissel (1977b). A
good, but dated, introduction to its application to quantum mechanics can be found in the book by Borowitz (1967) and to
solid mechanics in the review article by Steele (1976). What all this means is that the WKB approximation is probably a
very good idea since so many have rediscovered it in a wide variety of disciplines.
The WKB is a particular technique applied to a variety of problems, in particular linear differential equations of the
form:
ε2 y ′′ − q(x)y = 0 (508)
′′ 2
y + (λ p(x) − q(x))y = 0, λ≫1 (509)
′′ 2
y − q(εx) y = 0 (510)
It is not possible in general to obtain closed form solutions in terms of elementary functions to such differential equations
with variable coefficients.
From one point of view, the material in this section is also applicable to the equation:
εy ′′ + b(x, ε)y ′ + c(x, ε)y = 0 (511)

because the transformation
y = P (x)u, P (x) = exp(−B(x)/2ε), B ′ = b (512)

transforms the equation to:
ε2 u′′ − q(x)u = 0 (513)

with

b2 (x) b′ (x)
a(x, ε) = − −ε − c(x) . (514)
4 2
This change of variables is called the Sturm transformation (see Murdoch 1991). If one is interested in finding two
linearly independent solutions of (511), then this transformation is useful: it transforms the equation into (508) and the
solutions obtained using the techniques of this chapter can be transformed back into the original variable. But the Sturm
transformation drastically changes the nature of any initial or boundary conditions applied to this equation; in particular,
when boundary values are transformed from y to u, they may approach ∞ as ε → 0. Therefore the Sturm transformation is
almost entirely useless for asymptotic approximations to boundary value problems for small ε.
Unlike in multiple scales the secular terms are not explicitly suppressed, at least not at leading order. The underlying
additional condition is that the solution of the present type exists and that each higher order correction is no more secular
than its predecessor.
8.2 The method

In the same manner as was done for boundary layers and multiple scales, the ideas underlying the WKB method will be
developed by using it to solve a sample problem. The problem we begin with is the equation
110
ε2 y ′′ − q(x)y = 0. (515)
For the moment, we will not restrict the function q(x) other than to assume that it is smooth. Our intention is to
construct an approximation of the general solution of this equation. To motivate the approach that will be used, suppose
that the coefficient q is constant. In this case, the general solution of (515) with q constant (which is then second order,
linear, constant coefficient) is obtained via the exponential solution y = exp (αx) whence we obtain:
√
q
ε2 α2 − q = 0 =⇒ α = ± .
ε
Because the equation is second order and linear we know that the general solution comprises the superposition of two linearly
independent solutions and it follows that
1 √ 1 √
y(x) = a0 e− ε x q + b0 e ε x q . (516)
The hypothesis made in the WKB method is that the exponential solution in (516) can be generalized to
provide an approximate solution of (515). All that is necessary is to ensure the expansion is general enough so it can
handle the variable coefficient in the equation. In fact in many applications it is sufficient (and common) to seek a solution
in the form:
α
y ∼ eθ(x,ε)/ε ; θ(x, ε) = θ0 (x) + εα θ1 (x) + ... (517)
θ(x,ε)/εα
One should regard this substitution as a two stage process. One first makes the exact exponential substitution y ∼ e
to obtain a simplified equation; one then seeks a regular perturbation solution of the simplified problem by substituting
θ(x, ε) = θ0 (x) + εα θ1 (x) + . . .
Note Because the WKB method can only be applied to linear problems, there is a certainly flexibility possible in the
method. In particular the first task is usually to find (approximate) linearly independent solutions and form a general
solution by superposition. The boundary conditions are usually left till the end of the process and are often solved exactly
rather than approximately as has been the case using other asymptotic techniques (e.g. matched asymptotics, multiple
scales). This step is usually not thought provoking, but when using the WKB approximation we have to be a little more
careful. The reason is that when attempting to satisfy the boundary conditions it is very possible that the constants depend
on ε. It is therefore necessary to make sure this dependence does not interfere with the ordering assumed in the WKB
expansion in (517). See also section 8.5.2.
8.2.1 Examples
Example 1 Suppose q(x) = −e2x in (515).
Solution Let us use the substitution (517) so that:

′′ ′
θ(x, ε) dy 1 ′ d2 y θ θ2
y = exp ; = yθ ; 2 = y + 2 y.
ε dx ε dx ε ε
So the equation transforms (without approximation) to:
′′ ′
θ θ2 exp 2x
+ 2 + = 0 =⇒ (θ′ )2 + e2x + εθ′′ = 0.
ε ε ε2
Now seek a regular perturbation solution of this equation
′ ′ ′ ′′ ′′ ′′
θ = θ0 (x) + εθ1 (x) + O(ε2 ); θ = θ0 (x) + εθ1 (x) + O(ε2 ); θ = θ0 (x) + εθ1 (x) + O(ε2 )
and to O (ε) the equation becomes:
′ ′′ ′ ′
θ02 + exp 2x + ε(θ0 + 2θ0 θ1 ) + O(ε2 ) = 0.
At ε0 we have:
′ ′
θ02 + exp 2x = 0 ⇒ θ0 = ±i exp x ⇒ θ0 = ±i exp x + A
where A is a complex constant. At ε1 we have:
′′
′′ ′ ′ 1θ 1 ′1
1
ε : θ0 + 2θ0 θ1 = 0 ⇒ θ1 = − 0′ = − ⇒ θ1 = − x + B
2 θ0 2 2
so
111

θ0 ±i exp x + A 1 A ±i exp x 1
y = exp . exp θ1 = exp exp − x + B = exp B exp exp − x .
ε ε 2 ε ε 2
As the problem is linear, we can let A = B = 0 and we construct two linearly independent (complex valued) solutions:

i exp x 1 i exp x 1
y1 = exp exp − x ; y2 = exp − exp − x .
ε 2 ε 2
But exp x exp x
i exp x
exp = cos + i sin (518)
ε ε ε
and so we can construct two linearly independent (real) solutions:
x exp x x exp x
exp − cos ; exp − sin (519)
2 ε 2 ε
and the general solution is:
x exp x exp x
y = exp − α0 cos + β0 sin (520)
2 ε ε
as before.
The constants α0 , β0 are determined via the boundary conditions. For comparison, the exact solution is
y(x) = c0 J0 (λex ) + c1 Y0 (λex ), (521)
where
1 1
[bY0 (λ) − aY0 (λe)] and c1 = [aJ0 (λe) − bJ0 (λ)] .
c0 = (522)
d d
Here J0 and Y0 are Bessel functions and d = J0 (λe)Y0 (λ) − Y0 (λe)J0 (λ). The values obtained from the WKB approximation
(with boundary conditions y(0) = 1, y(1) = 0) and the exact solution are shown in Fig.48. It is clear that they are in very
good agreement.
Note We can also approach particular problems by using the general solution (574) for the case q(x) = −e2x . But it is
easier to construct an approximate solution for the particular example.
Note also that the ε1 problem could have been treated in the following way:
′′ Z Z
1 θ0 1 θ0′′ 1 dθ0′ 1
θ1′ = − ′ ⇒ θ1 = − , dx = − = − ln θ0′ + C.
2 θ0 2 θ0′ 2 θ0′ 2
Example 2 The WKB method is very useful for finding approximations of large eigenvalues. To illustrate how this
is done, consider the eigenvalue problem
y ′′ + λ2 e2x y = 0, for 0 < x < 1, (523)

where
y(0) = y(1) = 0, (524)

where we are interested in finding (large) eigenvalues λ and perhaps also the corresponding eigenvectors yλ (x). By definition
an eigenvalue is a value of λ for which (523) has a non-trivial solution for y(x). (Compare this with a matrix eigenvalue
problem where λ is chosen so that (A − λI)x = 0 has non-trivial x solutions.)
The assumption that λ is large leads precisely to the equation given in (515), where ε ≡ λ1 . The new feature of this
problem is that the parameter ε or λ is unknown (as is typical of eigenvalue problems). We are assuming that the problem
has (some) large eigenvalues. Using the WKB approximation given in (520), one determines the eigenvalues by using the
given boundary conditions. Specifically, using (520) with λ = 1/ε, the approximate solution for λ ≫ 1 is
e−x/2 [α0 cos (λex ) + β0 sin (λex )] . (525)

We comment at this point that the approach to dealing with boundary conditions in the WKB method is generally different
from previous approaches. Previously, the usual perturbation approach was to expand equation and boundary conditions
in terms of the small parameter hoping to obtain simplified versions of both at each order. Because the WKB method only
works for linear problems it is usually possible to write down a general solution via linearity and then to determine the
superposition constants by demanding that the approximate WKB solutions satisfy the full boundary conditions without
any further simplification. Unless otherwise noted that will be the approach here.
Applying the boundary conditions we see that:
112
α0 cos λ + β0 sin λ = 0; α0 cos λe + β0 sin λe = 0. (526)
It follows that
α0 = −β0 tan λ (527)

whence the second equation yields:
α0 cos λe + β0 sin λe = 0 ⇒ −β0 tan λ cos λe + β0 sin λe = 0 ⇒ − tan λ cos λe + sin λe = 0 (528)
and this eigenvalue equation can be rewritten as:
− sin λ cos λe + cos λ sin λe = 0 ⇒ sin λ(1 − e) = 0. (529)

Therefore, a first-term approximation of the large eigenvalues is λ(1 − e) = nπ i.e.
nπ
λ∼for integer n. (530)
e−1
The exact value of each eigenvalue can be shown to be determined from the denominator in (522), that is, the (exact)
eigenvalues are solutions of the equation
J0 (λe)Y0 (λ) − Y0 (λe)J0 (λ) = 0. (531)

To see just how well (530) does in approximating the eigenvalues, the relative error is plotted in Fig.49 for the first
twenty eigenvalues. Clearly, the approximation is very good and is much easier to use than (531). Moreover, the WKB
approximation does well even for the first few eigenvalues.
Example 3 Suppose that q(x) = x2 in (515) and consider the eigenvalue problem
x2
y ′′ + y; y(1) = 0, y(2) = 0 (532)
ε2
where the (large) eigenvalues 1/ε are to be determined. Using (517) for our WKB approximation we write:
θ(x,ε)
y=e ε , θ(x, ε) ∼ θ0 (x) + εθ1 (x) + . . .
and we obtain:
x2
ε0 θ0′2 = −x2 ⇒ θ0 = ±i +A (533)
2
1 θ0′′ 1 √
ε1 θ1′ = − =− ⇒ θ1 = − ln x + B (534)
2 θ0′ 2x
Thus two linearly independent solutions for y are:
i x2 i x2
i x2 √ eε 2 i x2 √ e− ε 2
exp − ln x ≡ √ ; exp − − ln x ≡ √ (535)
ε 2 x ε 2 x
and taking real parts a general solution is:
2 2
A cos x2ε + B sin x2ε
y= √ . (536)
x
Of course this is only an approximate general solution. One feature of the WKB method is that we can now treat this as if
it is an exact solution and satisfy the boundary conditions exactly by allowing the arbitrary constants to depend on ε. Using
the boundary conditions, we find that:
1
y(1) = 0 ⇒ A = −B tan (537)
2ε
and so
2 2
1
−B tan 2ε cos x2ε + B sin x2ε
y= √ (538)
x
Applying the other boundary condition we see that:
1 2 2 2 1
y(2) = 0 ⇒ − tan cos + sin = 0 ⇒ tan − tan =0 (539)
2 ε ε ε 2ε
113
1.0
Exact Solution
WKB Solution
0.5
Solution
0
-0.5
-1.0
0 0.2 0.4 0.6 0.8 1.0
x - axis
Figure 48: Exact solution and WKB solution given in (520) with ε = 0.1, y(0) = 1, y(1) = 0.
Using the identity:

tan A − tan B
tan(A − B) = (540)
1 + tan A tan B
provided the denominator 1 + tan A tan B is non-zero, we see that tan(A − B) = 0 when tan A − tan B = 0 and we find that
the eigenvalues are given by
2 1 1 2
tan( − ) = 0 ⇒ = nπ, n = 1, 2, 3 . . . (541)
ε 2ε ε 3
Example 4 Consider the equation:
d2 y
+ f (εt)y = 0
dt2
where ε ≪ 1 and f is a well behaved smooth function. εt is thus a slow time variable and the function f (εt) is slowly varying
in comparison with the t timescale. Let us define the problem in terms of the slow time-scale and write τ = εt, x(τ ) = y(t)
to obtain:
dy dx dτ dx d2 y 2
2d x
= =ε ; = ε
dt dτ dt dτ dt2 dτ 2
and the problem transforms to:
d2 x
ε2 + f (τ )x = 0
dτ 2
which is of WKB type and can be solved in the usual way.
8.2.2 General Discussion

The most distinguished characteristic of the WKB approximation is the assumed exponential dependence of the fast variation.
If we make this assumption, the first-term approximation of the solution is determined by solving two first-order differential
equations. The first one, the eikonal equation, is usually first-order and nonlinear, and it determines the fast variation in the
solution. It was not difficult to solve the eikonal equation for the example, but this can be one of the more difficult steps
in using the WKB method. The next equation to solve is the transport equation, and this determines the slow variation.
Fortunately, the transport equation is linear as are the equations for determining the higher-order terms in the expansion.
The WKB method is relatively easy to use and is capable of providing an approximation with less effort than other
methods. However, with this advantage comes limitations. The most significant is the requirement that the equation
be linear. This is needed to be able to cancel out the exponential term, as was done to obtain (568). For nonlinear problems,
the method of multiple scales can be used and is sometimes referred to as the nonlinear WKB (or nonlinear geometric optics)
method.
114
10-1
10-2
Relative error
10-3
10-4
10-5
0 5 10 15 20
n
Figure 49: Relative error, in absolute value, between the WB approximation (530) of the eigenvalues and the numerical
values.
8.3 Turning Points (WKB1-boundary layer-WKB2)

We will introduce the analysis for turning points using the example from the previous section. The equation to solve is
therefore
ε2 y ′′ − q(x)y = 0. (542)
As pointed out in deriving the WKB approximation, we must stay away from the points where q(x) is zero. To explain how
to deal with these points, we will initially consider the simplest case and then afterwards investigate how to extend the ideas
to more complex situations. So, to begin, suppose there is a single turning point xt where q(x) > 0 if x > xt and q(x) < 0 if
x < xt (see Fig.51). This means that the solution of (542) will be oscillatory if x < xt and exponential if x > xt .
The fact that the solution is oscillatory for negative q(x) can be understood if one considers the constant coefficient
equation y ′′ + λ2 y = 0, where λ ∈ R. The general solution in this case is y(x) = A0 cos(λx + φ0 ), which is clearly oscillatory.
A similar explanation can be given for the exponential solutions of y ′′ − λ2 y = 0 with general solution y = A exp(λx) +
B exp(−λx).
8.3.1 Example
As an example of a turning point problem, consider the following:
ε2 y ′′ = x(2 − x)y, for − 1 < x < 1, (543)

′
where y(−1) = y(1) = 1. In this case q(x) = x(2 − x), and so there is a simple turning point at x = 0 with q (0) = 2 (see
figure 8.3.1).
Using the substitution (517)
θ(x, ε)
y = exp ≡E
ε
we see that ! !
′ ′ ′
′ Eθ ′′ 1 Eθ 2
′′ E θ2
′′
y = ;y = Eθ + = θ +
ε ε ε ε ε
and the equation becomes:
′′ ′ ′′ ′
2 2
εEθ + Eθ = x(2 − x)E ⇒ εθ + θ = x(2 − x).
Now putting:
′ ′ ′
θ = θ0 (x) + εθ1 (x) . . . ; θ = θ0 (x) + εθ1 (x) . . .
the problem becomes:
115
Figure 50: Graph of x(2 − x).
′ ′ ′ ′′
θ02 + 2εθ0 θ1 + εθ0 = x(2 − x).
For −1 ≤ x < 0 i.e., x(2 − x) < 0 so write x(2 − x) = −x(x − 2) where x(x − 2) > 0. Then at ε0 we have
′ ′ p
θ02 = −x(x − 2) ⇒ θ0 = ±i x(x − 2)
so that

1 1 1 1
θ0 = ±i (x − 1)(x2 − 2x) 2 − ln(x − 1 + (x2 − 2x) 2 ) + C ≡ ±iΘ, say. (544)
2 2
At ε1 we then have:
′′ Z ′′ Z ′′ Z
′ θ0 1 θ0 1 θ0 dx ′ 1 1 ′ 1 ′
θ1 = − ′ ⇒ θ 1 = − ′ dx = − ′ ′ dθ 0 = − ′ dθ0 = − ln θ0 + A
2θ0 2 θ0 2 θ0 dθ0 2 θ0 2
1 p 1 p 1 1
= − ln(±i x(x − 2)) + A = − ln( x(x − 2)) − ln(±i) + A = − ln(x(x − 2)) + B
2 2 2 4
where B is a complex constant. Thus for −1 ≤ x < 0, the solution (in complex form) is

1 1
±i 21 (x − 1)(x2 − 2x) 2 − 12 ln(x − 1 + (x2 − 2x) 2 )
1
y ∼ exp exp − ln(x(x − 2)) + B .
ε 4
1
W.l.o.g. we can take B, C = 0 and exp − 41 ln(x(x − 2)) = (x(x − 2))− 4 . Using Euler’s result exp(ix) = cos x + i sin x we
thus see that a general solution is:

− 41 Θ Θ
yl ∼ (x(x − 2)) αl cos + βl sin . (545)
ε ε
For 0 < x ≤ 1 i.e., x(2 − x) > 0. Then at ε0 we have
′ ′ p
θ02 = x(2 − x) ⇒ θ0 = ± x(2 − x)
so that

1 1 1
θ0 = ± (x − 1)(2x − x2 ) 2 + arcsin(x − 1) + C ≡ ±Φ, say
2 2
At ε1 we then have:
116
Z ′′ Z ′′ Z
1 ′
′′
′ θ0 1 θ0 1 θ0 dx ′ 1 1 ′
θ1 = − ′ ⇒ θ1 = − ′ dx = − ′ ′ dθ0 = − ′ dθ0 = − ln θ0 + A
2θ0 2 θ0 2 θ0 dθ0 2 θ0 2
1 p 1
= − ln( x(2 − x)) + A = − ln(x(2 − x)) + B
2 4
where B is a complex constant. Thus for 0 < x ≤ 1, the solution (in complex form) is

±Φ 1
y ∼ exp exp − ln(x(2 − x)) + B .
ε 4
1
W.l.o.g. we can take B, C = 0 and exp − 14 ln(x(2 − x)) = (x(2 − x))− 4 . Thus we obtain the (real) general solution

1 Φ Φ
yR ∼ (x(2 − x))− 4 αR exp + βR exp(− ) (546)
ε ε
Solution near the turning point x = 0 Near x = 0 we rescale the problem by introducing the transition layer coordinate:
x y(x)
β
; Y (x) = γ x= (547)
ε ε
where β,γ will be determined from the analysis. Note that we have also rescaled Y . Even though the problem is linear,
matching later demonstrates that such a rescaling is necessary in order to perform the asymptotic matching. But note that
the fundamental balance in the equation is unaffected by the value of γ. The differential equation now becomes:
′′
ε2−2β Y (x) − εβ x(2 − εβ x)Y = 0.
For an asymptotic balance, we require 2 − 2β = β so β = 2/3. Now putting Y ∼ Y0 + . . . at leading order we find that:
′′
Y0 (x) − 2xY0 = 0 (−∞ < x < ∞) (548)
√
3
Putting s = 2x with Z(s) = Y0 (x) this transforms to the Airy equation (second order, linear, non-constant coefficient):
′′
Z (s) − sZ = 0 (−∞ < s < ∞)
with general solution in terms of Airy functions
√
3
√
3
Z(s) = aAi(s) + bBi(s) ⇒ Y0 (x) = aAi( 2x) + bBi( 2x). (549)
We note that we now have six undetermined constants. Two can be determined from the boundary conditions; the others
must be determined via asymptotic matching. We match the left hand solution with the transition layer and then the
transition layer with the right hand solution. We present the details of the matching for yl , Y0 . The rest is left as an exercise.
Matching of yl and Y0 Note that for small x < 0, the following holds:
√ √
1 2 1 1 2 1 2 2 3 2 2 3β
(x − 1)(x − 2x) − ln(x − 1 + (x − 2x) ) ∼
2 2 (−x) ≡
2 ε 2 (x)
2 2 3 3
so writing yl expanded for small x < 0 gives :
√ ! √ !!
1 2 2 3 2 2 3
1 1 αl cos (−x) 2 + βl cos (−x) 2
2 4 (−x) 4 3ε 3ε
in terms of the inner variable with x ≡ εβ x and expanding gives:

√ ! √ !!
1 2 2 3β 3 2 2 3β 3
(yl )inner = 1 β 1
αl cos ε 2 (−x) 2 + βl sin ε 2 (−x) 2 (550)
2 4 ε 4 (−x) 4 3ε 3ε
while the inner solution εγ Y0 (x) expanded for x → −∞ (see appendix) has the form:

εγ √ 2 3 π √ 2 3 π
√ 1 a cos 2 |x| −
2
+ b sin 2 |x| +
2
(551)
π |x| 4 3 4 3 4
Note: In terms of the outer variable (x = x/εβ ) the inner solution is

β
εγ− 4 √ 2 3β 3 π √ 2 3β 3 π
(εγ Y0 (x))outer = √ 1 a cos 2 ε 2 |x| 2 − + b sin 2 ε 2 |x| 2 +
π |x| 4 3 4 3 4
117
Comparing the expansions (550) and (551) it is clear that we must take
2 β 1
β= ;γ = − = − .
3 4 6
The relationships between the constants αl , βl , a, b then follow after some algebra. Inspection suggests that it is better to
rewrite (550) in the form αm cos(z + βm ) i.e., in the form:
√ !!
1 2 2 3β 3
(yl )inner = 1 β 1
αm cos ε 2 (−x) 2 + βm (552)
2 4 ε 4 (−x) 4 3ε
and now matching of (551) and (552) requires:
√
αm π 2
1 = a; βm = − ;b = 0 (553)
24 4
8.3.2 Examples
1. Suppose the problem is
ε2 y ′′ + sin(x)y = 0, for 0 < x < 2π,
where y(0) = a and y(2π) = b. This has three turning points: xt = 0, π, 2π. Because two of these are endpoints, there
d
will be two outer WKB approximations, one for 0 < x < π and one for π < x < 2π. Since dx sin(x) 6= 0 at x = π, then
the turning point xt = π is simple and is treated in much the same way as the one analysed above. The solution in the
transition layer at xt = 0 will be required to satisfy the boundary condition y(0) = a, and the solution in the layer at
xt = 2π will satisfy y(2π) = b.
2. The equation
ε2 y ′′ + p(x)y ′ + q(x)y = 0
differs from the one considered above because of the first derivative term. The WKB approximation of the general solution
is left as an exercise. Unlike before, the turning points now occur when p(x) = 0.
8.4 PDEs modelling wave propagation

The WKB method is quite useful for finding an asymptotic approximation of a traveling wave solution of a partial differential
equation. To illustrate this, we consider the problem

2 0<x<∞
uxx = µ (x)utt + α(x)ut + β(x)u, for , (554)
0<t
where
u(0, t) = cos(ωt). (555)

This is the equation for the displacement of a string that is damped (αut ) and has an elastic support (βu) . Because the
string is being forced periodically at the left end, the solution will develop into a wave that travels to the right. We wish to
try and construct such a solution without specifying initial conditions for the problem. (Initial conditions can be inferred
from the solution if we can construct it but here we are interested in seeing if the problem has a travelling wave solution).
We will, in fact, seek an asymptotic solution for the case when ω ≫ 1. We will motivate this further below.
To find such a solution, we will set u(x, t) = eiωt v(x) and then require that the function v(x) be consistent with the
observation that the wave moves to the right. Also, for the record, it is assumed that the functions in (554) are smooth with
µ(x) > 0 and α(x) and β(x) nonnegative.
Preamble to solution The equation is linear, but there is no obvious small parameter ε that can serve as the basis of
the WKB approximation. To motivate what will be done, suppose α = β = 0 and µ is constant. In this case, the problem is:

2 0<x<∞
uxx = µ utt , for ; u(0, t) = cos ωt, ω ≫ 1.
0<t
and we can construct plane wave (travelling periodic wave) solutions via the substitution
ω
u(x, t) = ei(ωt−kx) = eik( k t−x) , (556)
(where ω is the frequency, k is the wave number). To verify this, we substitute it into the constant coefficient equation
and find that k = ±ωµ i.e., u(x, t) = ei(ωt±ωµx) satisfies the PDE for all values of ω. This shows that for high frequencies
118

ω the waves have a relatively short wavelength (that is, λ = 2π k
≪ 1). This observation will serve as the basis of our
approximation, and in such a circumstance the small parameter is ε = ω1 . This assumption is what underlies the application
of the WKB method to many problems in optics, electromagnetics, and acoustics. For example, in the case of ultrasound,
the waves can have frequencies on the order of 106 Hz with wavelengths of about 3 × 10−2 m (in air). Even more extreme are
the waves for visible light, where the frequencies are in the neighborhood of 1015 Hz and the wavelengths are about 5 × 10−7
m. Finally note that (556) also demonstrates a travelling wave form (f (x − ct); c = ω/k).
′′
It may help in dealing with this problem to compare it with earlier ODE problems of the form ε2 y + q(x)y = 0. In the
present problem the x derivatives are relevant to the WKB method and the term µ2 (x)utt is playing the role of q(x)y in the
ODE problem.
WKB approach Based on the above discussion, we will construct an asymptotic approximation of the traveling wave
solution of (554) in the case of a (known) high frequency ω = 1ε , ε ≪ 1. This is done by generalising the form of (556) and
making the WKB assumption:

γ 1
u(x, t) ∼ eiωt v(x), v(x) = e−iω θ(x) u0 (x) + γ u1 (x) + · · · . (557)
ω
This may look more familiar if we write it in terms of ε as follows:
it −iθ(x)
u(x, t) ∼ e ε v(x), v(x) = e εγ [u0 (x) + εγ u1 (x) + · · · ]
This assumption is based on the fact that (554) has coefficients that depend on x but not on t. Thus the t dependent part
of our WKB substitution is unchanged. The x dependent part is adapted to (hopefully) have enough flexibility to deal with
the variable x dependent coeffficients. An alternative is of course to base the new approximation on (517) and write:
γ 1
u(x, t) ∼ eiωt v(x), v(x) = e−iω θ(x,ε)
; θ(x, ω) = θ0 (x) + θ1 (x) + . . . (558)
ωγ
or in terms of ε
it −iθ(x,ε)
u(x, t) ∼ e ε v(x), v(x) = e εγ ; θ(x, ε) = θ0 (x) + εγ θ1 (x) + . (559)
Solution We first make the substitution (558) i.e., u(x, t) ∼ eiωt v(x) in (554). Recalling that u(x, t) = v(x) exp(iωt) is
shorthand for u(x, t) = v(x) Re(exp iωt) = v(x) cos ωt = v(x) cos( εt ), we see that the problem reduces to:
′′
v = −ω 2 µ2 (x)v + iωα(x)v + β(x)v; v(0) = 1
1
and the first two terms suggest that this is a WKB type ODE with ε = ω ≪ 1. If we use a WKB approximation on this
ODE we now put
−iθ(x,ε)
v(x) = e ε
to obtain:
1 ′2 i ′′ 1 i ′2 ′′
− 2
θ − θ = − 2 µ2 (x) + α(x) + β(x) =⇒ −θ − iεθ = −µ2 (x) + iεα(x) + ε2 β(x).
ε ε ε ε
′2 ′ ′ ′
Now putting θ(x, ε) = θ0 (x) + εθ1 (x) + . . . we find that θ = θ02 (x) + 2εθ0 (x)θ1 (x) + O(ε2 ) and so
′2 ′ ′ ′′
−θ0 − 2εθ0 (x)θ1 (x) − iεθ0 = −µ2 (x) + iεα(x) + ε2 β(x) + O(ε2 ).
Thus ′
ε0 (θ0 )2 = µ2 (x).
Z x
′
Thus θ0 (x) = ±µ(x). The solutions of this eikonal equation are θ0 = ± µ(s)ds where we anticipate θ0 (0) = 0 will be
0
useful. At the next order we find that
′′ ′
′ ′ i θ0 + α(x) i µ (x) + α(x)
ε1 2θ0′ θ1 = −iθ0′′ − iα(x) =⇒ θ1 (x) = − ′ =− .
2 θ0 2 µ(x)
Thus
Z x Z x Z x
i α(s) i dµ i i α(s) i i
θ1 (x) = − ds − + ln a0 = − ds − ln(µ(x)) + ln a0 (560)
2 0 µ(s) 2 µ 2 2 0 µ(s) 2 2
for some arbitrary constant a0 . Putting the results together,
p imposing the boundary condition v(x = 0) = 1 i.e., θ(0, ε) ∼
θ0 (0) + εθ1 (0) = 0 (because v = exp θ/ε means that a0 = µ(0) as θ0 (0) = 0. Choosing the wave which propogates to the
right (cf. f (x ± ct)), we find that
119
s Z Z s Z Z x
µ(0) 1 x α(s) t 1 x µ(0) 1 x α(s)
u(x, t) ∼ exp − ds cos − µ(s)ds = exp − ds cos ωt − ω µ(s)ds .
µ(x) 2 0 µ(s) ε ε 0 µ(x) 2 0 µ(s) 0
(561)
Our solution is a traveling wave with an amplitude and phase that depend on x. The damping α(x), as might be expected,
causes the wave to be attenuated, while the coefficient µ(x) contributes both to the phase and amplitude. Interestingly, the
approximation is independent of the coefficient β(x) in (554). Its effect only appears at the next order (O(ε2 )).
8.5 WKB appendix

8.5.1 Other equivalent WKB substitutions
A slight generalisation of the WKB substitution already introduced is to seek a solution in the form:
α
y ∼ G(x, ε)eθ(x,ε)/ε ; G(x, ε) = y0 (x) + εα y1 (x) + ...; θ(x, ε) = θ0 (x) + εα θ1 (x) + .. (562)
where both G and θ are to be determined.
A similar approach when using the WKB method is to assume an asymptotic expansion of the form
α
y ∼ G(x, ε)eθ(x)/ε , G(x, ε) = y0 (x) + εα y1 (x) + · · · . (563)
where θ (a phase function) and y0 , y1... (amplitude functions) are to be determined. α will also be determined via asymptotic
balance. We emphasize that the three approaches are roughly equivalent; (517) is often a simpler way of proceeding from an
algebraic point of view.
One of the distinctive features of the WKB method is that it is fairly specific on how the solution depends on the fast
variation, namely, the dependence is assumed to be exponential. This limits the method, but it can make it easier to use
than either the multiple-scale or boundary-layer methods.
From (563) we get
α
y ′ ∼ (ε−α θ′ y0 + y0′ + θ′ y1 + · · · )eθ/ε (564)
and
α
y ′′ ∼ ε−2α θx2 y0 + ε−α θxx y0 + 2θx y0′ + θx2 y1 + · · · eθ/ε . (565)
These are sufficient for the problem we are considering. However, we will occasionally be considering third- and fourth-
order problems, and it is worth having the formulas in one place so they are easy to find when the time comes. So, for the
record, we also have
α
y ′′′ ∼ ε−3α θx3 y0 + ε−2α θx 3θx y0′ + 3θxx y0 + θx2 y1 + · · · eθ/ε (566)
and
α
y (4) ∼ ε−4α θx4 y0 + ε−3α θx2 6θxx y0 + 4θx y0′ + θx2 y1 + · · · eθ/ε . (567)
The next step is to substitute (564) and (565) into (515). Doing this, one finds that

1 1 ′′
ε2 (θ x )2
y 0 + θ y 0 + 2θ ′ ′
y 0 + (θ x )2
y 1 + · · · − q(x)(y0 + εα y1 + · · · ) = 0. (568)
ε2α εα
Something important happened in this step - the exponential dropped out. This would not have happened if the equation
were nonlinear. This point will be brought up again after the example has been completed. Now, balancing terms in (568),
one finds that α = 1, which leads to the following equation:
ε0 (θx )2 = q(x). (569)

This is called the eikonal equation ( from the theory of optics), and its solutions are
Z xp
θ(x) = ± q(s)ds. (570)
To obtain the first term in the expansion, we also need to find y0 (x), which is accomplished by looking at the O(ε) problem.
ε1 θ′′ y0 + 2θ′ y0′ + (θx )2 y1 = q(x)y1 .
120
This is the transport equation. Since θ(x) satisfies the eikonal equation, it follows that the last two terms are just 0 and the
equation reduces to θ′′ y0 + 2θ′ y0′ = 0. This is a first order linear variables separable ODE for y0 (x) where the coefficient
function θ(x) is a known function given in (570). The equation can be written as
Z s=x ′′ Z s=x
dy0 θ′′ (x) dy0 θ′′ (x) θ (s) 1 1
=− ′ y0 ⇒ =− ′ dx ⇒ ln y0 = − ds + A = − dθ′ + A (571)
dx 2θ (x) y0 2θ (x) 2θ′ (s) 2 θ′ (s)
1 1 p ln y
= − ln θ′ (x) + A = − ln( q(x)) + A ⇒ p 0 =A (572)
2 2 ln( 4 q(x))
whence we obtain
c
y0 (x) = p , (573)
4
q(x)
where c(= exp(A)) is an arbitrary constant.
We have therefore found that a first-term approximation of the general solution of (515) is
 Z x Z x 
√ √
− 1ε q(s)ds 1
q(s)ds
 ε 
y ∼ q(x)−1/4 a0 e + b0 e , (574)
where a0 and b0 are arbitrary, possibly complex, constants. It is apparent from this that the function q(x) must be nonzero.
The values of x where q(x) is zero are called turning points and are discussed elsewhere in the chapter.
8.5.2 Second Term of a WKB Expansion

To obtain a measure of the error in using (574) to approximate the solution, we will calculate the second term in the
expansion. This is not difficult; carrying out the calculation determines that the ε2 equation is
ε2 : θ′′ y1 + 2θ′ y1′ + (θx )2 y2 + y0′′ = q(x)y2 .

Using the eikonal equation (569) reduces this to
2θ′ y1′ + θ′′ y1 = −y0′′ .

In principle, this linear equation can be solved (exactly) using variation of parameters. However, letting y1 (x) = y0 (x)w(x),
one finds from the previous equation that
Z
1 θxx 1 x (θxx )2
w(x) = k + + dx, (575)
4 (θx )2 8 (θx )3
where k is an arbitrary constant.
Now, given a small value of ε, the expansion in (574) will be an accurate approximation of the solution if εy1 (x) is much
smaller than y0 (x). This can be guaranteed if |εw(x)| ≪ 1. From the formula for w(x) given in (575), we can express this
requirement in terms of q(x) and its first derivative. In particular, for x0 ≤ x ≤ x1 , we will have an accurate approximation
if the following holds:
Z x1 ′
1 q ′ q
dx
ε |k| + 3/2 4 + q ≪ 1, (576)
32 q ∞ x0
where |h(x)|∞ ≡ maxx0 ≤ x ≤ x1 |h(x)| . As might be expected, this shows that the expansion becomes nonuniform as a
turning point is approached, that is, when one gets close to a point where q(x) = 0. For a more extensive investigation into
the theory underlying the WKB method, readers should consult the book by Olver (1974).
The general solutions given in (574) and (575) contain constants that are determined from boundary conditions. This
step is usually not thought provoking, but when using the WKB approximation we have to be a little more careful. The
reason is that when attempting to satisfy the boundary conditions it is very possible that the constants depend on ε. It is
therefore necessary to make sure this dependence does not interfere with the ordering assumed in the WKB expansion in
(563). To illustrate this situation, consider the boundary-value problem
ε2 y ′′ + e2kx y = 0, for 0 < x < 1. (577)

This is a generalization of the examples considered earlier (where k = 1).
Since θx = ±iekx , we get, from (575),
ik −kx
w(x) = k ∓ e . (578)
8
121
q
x(t) x
′
Figure 51: In the analysis of the turning point, q(x) is assumed to have a simple zero at x = xt with q (xt ) > 0. This will
enable us to use a linear approximation for q(x) near the turning point.
A two-term expansion of the general solution of (577) is therefore

i kx ik
y ∼ e−kx/2 a0 e− kε e 1 + ε a1 − e−kx (579)
8

i kx ik
+ b0 e kε e 1 + ε b1 + e−kx .
8
One can use this to impose the boundary conditions. As in the example shown in Fig.48, we take k = 1, y(0) = 1, and
y(1) = 0. One finds in this case that in (579)
cos λ(e − 1) − e−1−i(e−1)/ε

a1 = − ,
8 sin λ(e − 1)
where λ = 1ε . One also finds that b1 is the complex conjugate of a1 . This shows that for this example the second term in the
combined WKB expansion in (579) is small except at, or near, the resonant points where λ(e − 1) = nπ. In particular, if
ε = 0.1, then one finds that the maximum value of a1 − ik 8 e
−kx
for 0 ≤ x ≤ 1 is about 0.08. This explains why the curves
in Fig.48 are essentially indistinguishable from each other.
8.5.3 Turning point problems: the general case

In the general case (542), we can use the WKB approximation on either side of the turning point. This gives the following
approximation for the general solution:

yL (x), for x < xt
y∼ , (580)
yR (x), for xt < x
where
 Z x Z x 
√ √
− ε1 q(s)ds 1
q(s)ds
 ε

yR = q(x)−1/4 aR e xt + bR e xt  (581)
and
 Z xt Z xt 
√ √
− 1ε q(s)ds 1
q(s)ds
 ε 
yL = q(x)−1/4 aL e x + bL e x  (582)
These expressions come directly from (574) except that we have now fixed one of the endpoints in the integrals at the
turning point. This particular choice is not mandatory, but it does simplify the formulas in the calculations to follow. It is
important to recognize that the coefficients in (581) and (582) are not independent ; we must find out how they are connected
by investigating what is happening in a transition layer centered at x = xt . This is very similar to what happened when we
were studying interior-layer problems in the chapter on matched asymptotic expansions. When we finish, the approximation
of the general solution will only contain two arbitrary constants rather than four as in (580).
122
8.5.4 Solution in Transition Layer
To determine the solution near the turning point, we introduce the transition-layer coordinate
x − xt
x= . (583)
εβ
We know the point xt , but we will have to determine the value for β from the analysis to follow. Now, to transform the
differential equation, we first use Taylor’s theorem to conclude

q xt + εβ x ∼ q (xt ) + εβ xq ′ (xt ) + · · · . (584)
′
We will assume q(x) has a simple zero at xt , and so q (xt ) 6= 0. With this, and letting Y (x) denote the solution in this layer,
we get from (542) that

ε2−2β Y ′′ − εβ xqt′ + · · · Y = 0, (585)
where qt′ ≡ q ′ (xt ) is a known constant. For balancing, we need 2 − 2β = β, and thus β = 23 .
The appropriate expansion of the solution in this region is
Y ∼ εγ Y0 (x) + · · · . (586)
Introducing this into (585), we get the following equation to solve:
Y0′′ − xqt′ Y0 = 0, for − ∞ < x < ∞. (587)

This can be transformed into an equation whose solutions are known. Letting s = (qt′ )1/3 x, we get Airy’s equation, which is
d2
Y0 − sY0 = 0, for − ∞ < s < ∞. (588)
ds2
Because of its importance in applied mathematics, Airy’s equation has been studied extensively. It can be solved using power
series expansions or the Laplace transform. One finds that the general solution can be written as
Y0 = aAi(s) + bBi(s),
where Ai(·) and Bi(·) are Airy functions of the first and second kinds, respectively. The definitions and some of the properties
of these functions are given in an Appendix. We are now able to write the general solution of the transition-layer equation
(587) as
h i h i
Y0 (x) = aAi (qt′ )1/3 x + bBi (qt′ )1/3 x . (589)
From (581), (582), and (589) we have six undetermined constants. However, the solution in (589) must match with the
outer solutions in (581), (582). This will lead to connection formulas between the constants, and these will result in two
arbitrary constants in the general solution.
8.5.5 Matching
Using intermediate variable The solution in the transition region must match with the outer solutions given in (581)
and (582). In what follows, keep in mind that qt′ > 0. To do the matching, we will use the intermediate variable
x − xt
xη = , (590)
εη
where 0 < η < 2/3. Before matching the solutions, note that the terms in the two outer solutions contain the following: For
x > xt
Z xp Z xt +εη xη p Z xt +εη xη p
q(s)ds = q(s)ds ∼ (s − xt )qt′ ds (591)
xt xt xt
2 3/2
= εr
3
and
−1/4
q(x)−1/4 ∼ [qt + (x − xt ) qt′ ] (592)
−1/6
=ε (qt′ )−1/6 r−1/4 ,
where r = (qt′ )1/3 εη−2/3 xη .
123
Matching for x > xt : Using the asymptotic expansions for the Airy functions given in an Appendix, one finds that

Y ∼ εγ Y0 εη−2/3 xη + · · · (593)
γ γ
aε 2 3/2 bε 2 3/2
∼ √ 1/4 e− 3 r + √ 1/4 e 3 r ,
2 πr πr
and for the WKB solution in (581),
ε−1/6 − 23 r 3/2 2 3/2

yR ∼ a e + b e 3r . (594)
R R
(qt′ )1/6 r1/4
For these expressions to match, we must have γ = − 16 ,
a 1/6 b 1/6
aR = √ (qt′ ) and bR = √ (qt′ ) . (595)
2 π π
Matching for x < xt : The difference with this case is that xη < 0, which introduces complex numbers into (582). As
before, using the expansions for the Airy functions given in the Appendix,

Y ∼ εγ Y0 εη−2/3 xη (596)

aεγ 2 3/2 π bεγ 2 3/2 π
∼ √ 1/4
cos |r| − +√ 1/4
cos |r| +
π |r| 3 4 π |r| 3 4
ε γ
−iπ/4 iπ/4 i 32 |r|3/2 iπ/4 −iπ/4 −i 32 |r|3/2
= √ 1/4
ae + be e + ae + be e .
2 π |r|

In the last step leading to (596), the identity cos θ = 21 eiθ + e−iθ was used. As for the WKB expansion, from (582),
ε−1/6 2 3/2 3/2

−π 2
+π
yL ∼ aL ei( 3 |r| 4) + bL e−i( 3 |r| 4) . (597)
(qt′ )1/6 |r|1/4
Matching (596) and (597) yields the following:
(qt′ )1/6
aL = √ (a + ib) (598)
2 π
and
(qt′ )1/6
√ (ia + b) = iaL ,
bL = (599)
2 π
where aL is the complex conjugate of aL . Equations (595), (598), and (599) are known as connection formulas. They
constitute a system of equations that enable us to solve for four of the constants in terms of the remaining two.
van Dyke matching The solution in the transition region must match with the outer solutions given in (581) and (582).
In what follows, keep in mind that the known constant qt′ > 0. To do the matching, recall that the outer solutions are known
to O(1) and the inner to O(εγ ) and
x − xt
x= ; x = xt + ε2/3 x. (600)
ε2/3
In the left outer solution x < xt while in the right outer x > xt . Before matching the solutions, note that the following results
pertaining to the outer solutions: for x > xt and x − xt ≪ 1
q(x) ∼ [qt + (x − xt )qt′ + . . .] ∼ 0 + (x − xt )qt′ (601)

−1/4 −1/4
q(x) ∼ [qt + (x − xt ) qt′ ] =ε −1/6
(qt′ )−1/6 r−1/4 ,
Z xp Z xt +εη xη p Z xt +εη xη p
q(s)ds = q(s)ds ∼ (s − xt )qt′ ds (602)
xt xt xt
2 3/2
= εr
3
where r = (qt′ )1/3 x.
124
Matching for x > xt : The inner solution is available to O(εγ ) where γ is to be determined in the matching. The
outer is available to O(1). Using the asymptotic expansions for the Airy functions given in the Appendix for large argument
(x = x−x
ε2/3
t
), one finds that
(Y )i ∼ εγ Y0 (x) + · · · (603)
aεγ 2 3/2 bεγ 2 3/2
∼ √ 1/4 e− 3 r + √ 1/4 e 3 r ,
2 πr πr
and for the WKB solution in (581),
ε−1/6 − 23 r 3/2 2 3/2

(yR )o ∼ a R e + b R e 3r . (604)
(qt′ )1/6 r1/4
For these expressions to match, we must choose γ = − 16 ,
a 1/6 b 1/6
aR = √ (qt′ ) and bR = √ (qt′ ) . (605)
2 π π
Matching for x < xt : The difference with this case is that x < 0, which introduces complex numbers into (582). As
before, using the expansions for the Airy functions given in the Appendix,
(Y )i ∼ εγ Y0 (x) (606)
γ
γ

aε 2 3/2 π bε 2 3/2 π
∼ √ 1/4
cos |r| − +√ 1/4
cos |r| +
π |r| 3 4 π |r| 3 4
ε γ
−iπ/4 iπ/4 i 32 |r|3/2 iπ/4 −iπ/4 −i 32 |r|3/2
= √ 1/4
ae + be e + ae + be e .
2 π |r|

In the last step leading to (596), the identity cos θ = 21 eiθ + e−iθ was used. As for the WKB expansion, from (582),
ε−1/6 2 3/2 3/2

−π 2
+π
(yL )o ∼ 1/4
aL ei( 3 |r| 4) + bL e−i( 3 |r| 4) . (607)
(qt′ )1/6 |r|
Matching (596) and (597) yields the following:
(qt′ )1/6
aL = √ (a + ib) (608)
2 π
and
(qt′ )1/6
√ (ia + b) = iaL ,
bL = (609)
2 π
where aL is the complex conjugate of aL . Equations (595), (598), and (599) are known as connection formulas. They
constitute a system of equations that enable us to solve for four of the constants in terms of the remaining two.
8.5.6 Summary
Solving the connection formulas in (605), (608) and (609) shows that the WKB approximation of the solution of (542) is
( 1

|q(x)|1/4
2aR cos 1ε θ(x) − π4 + bR cos 1ε θ(x) + π4 , for x < xt
y(x) ∼ 1 1
(610)
1
q(x)1/4
aR e− ε k(x) + bR e ε k(x) , for xt < x,
where
Z xt p
θ(x) = |q(s)|ds (611)
x
and
Z xp
k(s) = q(s)ds. (612)
xt
It should be remembered that this expansion was derived under the assumption that x = xt is a simple turning point with
q ′ (xt ) > 0. The accuracy of this approximation near the turning point depends on the specific problem. However, one can
show that in general one must require ε2/3 ≪ |x − xt | (Exercise). Also, as expected, we have ended up with an expansion for
the solution of (542) that contains two arbitrary constants (aR , bR ) which would be determined from the boundary conditions.
125
6 yL
yR
4
Y
2
0
-2
a
4
2
0
yL Y
-2
yR Numerical
4
b -1.0 -0.5 0 0.5 1.0
x - axis
Figure 52: (a) The transition layer and WKB approximations of (586) and (610). (b) A comparison of the transition layer
and WKB approximations of (586) and (610) and the numerical solution of (613). ε = 0.025.
8.5.7 Example
We reconsider the following problem (already solved above) and put it into the context of the general case:
ε2 y ′′ = x(2 − x)y, for − 1 < x < 1, (613)

′
where y(−1) = y(1) = 1. In this case q(x) = x(2 − x), and so there is a simple turning point at x = 0 with q (0) = 2. The
WKB approximation in (610) therefore applies, where, from (611) and (612),
1 1 h i
θ(x) = (1 − x)(x2 − 2x)1/2 − ln 1 − x + (x2 − 2x)1/2
2 2
and
1 1 π
k(x) = (x − 1)(2x − x2 )1/2 + arcsin(x − 1) + .
2 2 4
The constants aR and bR are found from the two boundary conditions. For example, since q(1) = 1 and k(1) = π/4, then the
π π
condition y(1) = 1 leads to the equation aR exp(− 4ε ) + bR exp( 4ε ) = 1. The resulting approximation obtained from (610) is
shown in Fig. 52 along with the approximation from the transition layer. The singularity in yR and yL at x = 0 is evident
in this plot as is the matching by the Airy functions from the transition layer. To demonstrate the accuracy of the WKB
approximation even when a turning point is present, in Fig.52 the numerical solution is plotted along with the curves from
Fig.52 . Within their respective domains of applicability, the WBK approximations are essentially indistinguishable from
the numerical solution.
One can use (610) in a straightforward manner to determine the expansion when x = xt is a simple turning point with
q ′ (xt ) < 0. Carrying out the details, it is found in this case that (see Exercise 1)
(
1 − 1ε k(x) 1
ε k(x)
q(x) 1/4 αR e + β R e , for x < xt
y(x) ∼ 1 1 ε
1 π
(614)
|q(x)|1/4
2α R cos ε θ(x) + 4 + β R cos ε θ(x) − 4 , for xt < x,
where θ(x) and k(x) are given in (611) and (612). Again, we have a solution that contains two arbitrary constants (αR , βR ) .
We have investigated the case of a single first-order turning point located at x = xt . The solution is in three pieces: the
two WKB approximations (yL , yR ) on either side of the turning point and the transition-layer solution in between. The
first-term approximation of the solution in the transition layer was found to be a solution of Airy’s equation (587), which
is the prototype equation for a simple turning point. In fact, it is possible to find a composite first-term expansion of the
solution in terms of Airy functions. This was first derived by Langer (1931); the result is
" 2/3 ! 2/3 !#
k(x)1/6 3k 3k
y(x) ∼ 1/4
a0 Ai + b 0 Bi , (615)
|q(x)| 2ε 2ε
where k(x) is given in (612). The derivation of (615) is left as an exercise. Also, in connection with the WKB approximations
given in (610),
126
1/6 1/6
a0 2ε b0 2ε
bR = √ and aR = √ .
2 π 3 π 3
Because of the value of having a composite expansion, there have been numerous generalizations of Langer’s result. A
discussion of some of these can be found in Nayfeh (1973).
In referring to the turning point as simple, it is meant that q(x) has a simple zero at x = xt (i.e., q(xt ) = 0 but q ′ (xt ) 6= 0).
Higher-order turning points do arise; for example, a second-order turning point, at x = 0, is
ε2 y ′′ − x2 ex y = 0. (616)
′ ′′
The reason this is second order is simply that q(xt ) = q (xt ) = 0 but q (xt ) 6= 0. It is also possible to have one of
fractional order (e.g., when q(x) = x1/3 ex ). The prototype equations in these cases are discussed in Holmes (1995). Other
types do occur, such as logarithmic, although they generally are harder to analyze.
Another complication that can arise is that the position of the turning point may depend on ε. This can lead to coalescing
turning points; an example of this is found in the equation
ε2 y ′′ − (x − ε)(x + ε)y = 0. (617)

This has simple turning points at xt = ±ε, but to the first-term they look like one of the second order at x = 0. Situations
such as this are discussed in Steele (1976).
8.5.8 Energy Methods

To connect what we have done with the energy methods that are also used to study wave problems, we first need to determine
an energy equation associated with (554). One way to do this is to multiply (554) by the velocity ut (Exercise). Doing this
and collecting terms, one obtains an equation of the form
∂t E + ∂x S = −Φ, (618)
where E = 21 µ2 (ut )2 + 12 (ux )2 + 21 βu2 , S = −ut ux , and Φ = α(ut )2 . As is always the case, the problem we are considering has
been nondimensionalized. However, if we were to convert back to dimensional coordinates, we would find that E corresponds
to the sum of the kinetic and potential energy densities, S is the energy flux, and Φ = α(ut )2 is the rate of energy dissipation
(or, simply, the dissipation function). We will refer to E, S, and Φ in these physical terms, although it is understood that
they are actually nondimensional.
Wave solutions are of particular importance in most areas of application. Consequently, numerous methods have been
devised to find approximations of them. Of interest here are those based on using averages of the energy. The assumption
made is that for a high frequency
u(x, t) ∼ A(x) cos [ωt − ϕ(x)] , (619)

where the amplitude A is slowly varying in comparison to the change in the phase ϕ (i.e., the wave is slowly modulated). It
should be noted that the phase velocity vph for this wave is
ω
vph ≡ . (620)
ϕx
With this, curves x = x(t) in the (x, t)−plane that satisfy x′ = vph are called phase lines. To help in understanding the next
result, recall that, from the WKB expansion, the wavelength of the fast variation is O(ω −1 ).
Claim 1: The eikonal equation for (554) is equivalent to the requirement that the total energy remain constant (to the
first term) between any two phase lines that are O(ω −1 ) apart.
The total energy Etot (t) between two phase lines is
Z x2 (t)
Etot = E(x, t)dx, (621)
x1 (t)
where x′1 = x′2 = vph (and x1 < x2 ). With the assumption in (619), we have
1 2 2 2
E∼ A µ ω + ϕ2x [1 − cos 2 (ωt − ϕ)] , (622)
4
1
S∼ ωϕx A2 [1 − cos 2 (ωt − ϕ)] , (623)
2
and
1 2 2
Φ∼ αω A [1 − cos 2 (ωt − ϕ)] . (624)
2
127
Now, using the energy equation and the fact that x1 − x2 = O(ω −1 ) gives us
Z x2 (t)
d
Etot = E2 x′2 − E1 x′1 + ∂t E(x, t)dx
dt x1 (t)
∼ E2 x′2 − S2 − E1 x′1 + S1 ,
where Ei = E(xi , t) and Si = S(xi , t) for i = 1, 2. Using the expansions in (622) - (624) yields
ωA2 2 2
Ex′ − S ∼ ω µ − ϕ2x sin2 (ωt − ϕ) .
2ϕx
d
Recall that the eikonal equation for this problem is ω 2 µ2 = ϕ2x . From this it is clear that the equation dt Etot = 0, to the
first term, is equivalent to the eikonal equation.
Claim 2: The transport equation corresponds to the balance of energy over a period.
Given a function g(t), the average to be used is
Z 2π/ω
ω
gavg = g(t)dt.
2π 0
Averaging the energy equation (618) and using (622), one obtains
∂x Savg = −Φavg . (625)

Using (623) and (624) yields
∂x (ϕx A2 ) = −αA2 .
The solution of this is the same as the solution of the transport equation.
The above results are interesting because they give a physical interpretation of the WKB approximation. However, these
claims are not necessarily true for other problems, and so, one must use caution in trying to use them to construct the WKB
approximation. It is much better to rely on the original expansion. Also, the above discussion involved the phase velocity,
but the group velocity plays an even more fundamental role in energy propagation. The interested reader should consult
Whitman (1974) and Lighthill (1965).
8.5.9 Slender Body Approximations

Another important situation where the WKB approximation can be used to obtain an asymptotic approximation of traveling
wave solutions arises when the geometry is long and thin. This sort of situation is referred to as slender body theory.
Example To illustrate this, consider the following problem:

2 2 0<x<∞
ε uxx + uyy = µ (x)utt , for , (626)
−G(x) < y < G(x)
where
u(x, y, t) = 0 for y = ±G(x), (627)

and
u(0, y, t) = f (y) cos(ωt). (628)

This problem corresponds to finding the vertical displacement u(x, y, t) of an elastic membrance that is much longer than
it is wide. The membrane is fixed along its lateral sides (y = ±G) and is forced periodically at the left end (x = 0). In effect,
this forcing is generating waves at x = 0 that propagate down the membrane. We are interested in finding these waves and
determining their basic properties. For example, once a wave is sent down the membrane, do we ever see it again? You might
expect that because the membrane is infinitely long that it just keeps traveling away from the end at x = 0. However, as will
be seen later, it is actually possible for the wave to travel a certain distance along the membrane and then turn around and
come back to x = 0. It is not obvious why this should happen without help from the analysis to follow. In this example, the
forcing frequency ω is not restricted, other than being positive.
To find the traveling wave solutions for small ε, we use the WKB approximation
α
u(x, y, t) ∼ ei[ωt−θ(x)/ε ] [u0 (x, y) + εα u1 (x, y) + · · · ] . (629)
128

ε2−2α −θx2 u0 − iεα θxx u0 + 2θx ∂x u0 + θx2 u1 + · · · + ∂y2 u0 + εα ∂y2 u1 + · · · (630)
2 2 α
= −ω µ (u0 + ε u1 + · · · ) .
Balancing the terms in this equation, it follows that α = 1. This leads to the following problems:

O(1) ∂y2 u0 + ω 2 µ2 − θx2 u0 = 0,
u0 = 0 for y = ±G(x).
This is the eikonal problem. It is actually an eigenvalue problem where θx2 is the eigenvalue and u0 is an eigenfunction.
1/2
Setting λ = ω 2 µ2 − θx2 , the solution of this can be written as
u0 (x, y) = A(x) sin [λ(y + G)] , (631)

where to satisfy the boundary conditions we must have 2λG = nπ. Therefore,
p
θx = ± ω 2 µ2 − λ2n for n = 1, 2, 3, . . . , (632)
where λn (x) = nπ/(2G). This means that the function θx depends on the mode number (n) as well as on the position (x).
Also, as long as θx is real and nonzero, the wave propagates along the x−axis. However, the higher modes (where n is large
enough that λn > ωµ) do not propagate since the wavenumber is imaginary.
O(ε) ∂y2 u1 + λ2n u1 = i(θxx u0 + 2θx ∂x u0 )

u1 = 0 for y = ±G(x)
This is the transport problem and we will use it to complete the determination of u0 . Because it is an inhomogeneous version
of the eikonal problem, we will have to impose a condition on the right-hand side of the transport equation to guarantee a
solution. This situation is covered by the Fredholm Alternative Theorem. To obtain the condition, we multiply the equation
by u0 and then integrate with respect to y. Doing this, one obtains the following:
Z G Z G

u0 ∂y2 u1 + λ2n u1 dy = i u0 (θxx u0 + 2θx ∂x u0 ) dy.
−G −G
Integrating by parts twice with the left integral, and using the eikonal problem, shows that the left-hand side is zero. Thus,
Z G
0= ∂x θx u20 dy. (633)
−G
This equation for u0 is known as a solvability condition, and it must be satisfied for the transport problem to have a solution.
To solve (633), we use Leibnitz’s rule to obtain
Z G
d
0= θx u20 dy − G′ θx u20 Gy=−G .
dx −G
Since u0 = 0 for y = ±G, we have

Z G
0x u20 dy = a,
−G
where a is a constant. Thus, from (631),

a
A= p .
θx G(x)
Putting the results together, we have found that the WKB approximations for the modes have the form
a
u(x, y, t) ∼ p ei[ωt±θ(x)/ε] sin [λn (y + G)] , (634)
θx G(x)
where
Z xp
θ(x) = ω 2 µ2 − λ2n dx (635)
and
129
nπ
λn = , for n = 1, 2, 3, . . . . (636)
2G(x)
It is apparent from this that not every mode is going to produce a traveling wave. This is because θ(x) is imaginary whenever
λn > ωµ. For example, as mentioned earlier, the higher modes (where λn is large) do not propagate. In any case, to obtain the
solution of the original problem, these modes should be added together and the constants an then determined by satisfying
(628).
Before starting the analysis for the transition-layer regions, it is worth making a few observations. First, the eikonal
problem differs significantly from the one obtained in (8.4). However, even though it is a linear boundary-value problem,
the equation that determines θx is nonlinear. In other words, it is still true that the equation that determines θ(x) is a
non-linear, first-order differential equation. Second, even though the next-order problem differs from what we saw earlier, it
is still true that the transport equation, which we referred to as a solvability condition for the problem, is a linear, first-order
linear differential equation.
Solution in Transition Region The turning points in the WKB approximation (634) occur when θx = 0. Using (632),
these points satisfy
nπ
µ(xt ) = . (637)
2ωG(xt )
We will consider the case when there is a single turning point 0 < xt < ∞. To guarantee there is only one, it is assumed
d
that dx θx2 < 0 (i.e., θx2 is a decreasing function of x, so if it does go through zero it does not do so again).The first step in
the procedure to find the solution in the transition region is to introduce the transition-layer coordinate
x − xt
x= . (638)
εβ
Letting U (x, y, t) denote the solution in this region, we assume that
U (x, y, t) ∼ εγ eiωt {U0 (x) sin [λn (y + Gt )] + · · · } , (639)

where Gt = G(xt ). Substituting (638) and (639) into the differential equation (626) produces the following result:
ε2−2β U0′′ + θx2 U0 + · · · = 0, (640)

where
d
θx2 xt + εβ x ∼ θx2 (xt ) + εβ x θx2 (xt ) + · · · (641)
dx
β d 2
∼ ε x θx (xt ) + · · · .
dx
For the terms to balance in (640), we need 2 − 2β = β, and thus β = 2/3.
The equation to be solved is thus
U0′′ = kxU0 , for − ∞ < x < ∞, (642)

where

d µ′t G′
k = − θx2 (xt ) = −2ω 2 µ2t + t > 0. (643)
dx µt Gt
This is an Airy equation, and the general solution is

U0 = aAi k 1/3 x + bBi k 1/3 x . (644)
Matching What remains is the matching. We can use the WKB approximation of the nth mode on either side of the
turning point. So,

uL (x, y, t), for 0 ≤ x < xt
u (x, y, t) ∼ , (645)
uR (x, y, t), for xt < x < ∞
where
1 n o
uL = p aL ei[ωt−θ(x)/ε] + bL ei[ωt+θ(x)/ε] sin [λn (y + G)] (646)
θx G(x)
130
y
y = G(x)
xt x
y = -G(x)
Figure 53: Tapered membrane which leads to the WKB approximation given in (652).
and
1
uR = p aR eiωt−ϕ/ε sin [λn (y + G)] . (647)
|θx | G(x)
Also,
Z xt p Z xp
θ= ω 2 µ2 − λ2n ds and ϕ = λ2n − ω 2 µ2 ds. (648)
x xt
The endpoints in the above integrals are chosen to make the calculations in the matching somewhat simpler.
Before discussing the matching, note that for the solution to be bounded we must have b = 0 in (644). It is for this same
reason that the bR term is not included in (647). Now, the matching is very similar to what was done in (8.2) to derive (598)
and (599). One finds that in (639), γ = −1/6. Using the asymptotic properties of the Airy functions given in Appendix A1,
the connection formulas are found to be
√
2 π
a= 1/2
e−iπ/4 aL , (649)
k 1/6 Gt
bL = −iaL , (650)
and
1/2
k 1/6 Gt
aR = √ a = e−iπ/4 aL . (651)
2 π
Another example To see more clearly what we have derived, suppose the membrane is tapered and has a uniform density
so µ is constant (see Fig.53). Also, suppose the frequency is high enough that the nth mode has a turning point that satisfies
0 < xt < ∞. The WKB solution for this mode is
αn
u(x, y, t) ∼ p Vn (x) sin [λn (y + G)] cos (ωt) , (652)
θx G(x)
where
1 π

sin ε θ(x) + 4 for 0 ≤ x < xt
Vn (x) = −ϕ(x)/ε . (653)
e , for xt < x < ∞
The constant αn in (652) is determined from boundary condition (628), and the functions θ and ϕ are given in (648).
We are now in a position to summarize the situation. From (646) it is seen that in the region 0 ≤ x < xt the solution for
the mode consists of the superposition of left and right traveling waves. As the right running wave moves toward x = xt , its
amplitude grows since both G and θx are decreasing. Similarly, the wave that is reflected from the transition region decreases
in amplitude as it travels away from the turning point. Together they result in a standing wave, as described in (652) and
131
(653). The portion of the disturbance that gets through the transition region ends up decaying exponentially, as shown in
(653). This is known as the evanescent region.
There are several interesting observations that can be made using our solution. First, if the frequency is decreased, then
the position of the turning point, which is a solution of the equation
nπ
ωµ = ,
2G(xt )
moves to the left. Suppose ω = ωc , which has the turning point xt = 0. In this case ωc corresponds to a frequency below
which waves do not propagate without decaying exponentially. This is called a cutoff frequency. In this example, it is referred
to as a low-frequency cutoff because frequencies below this value do not propagate. The second observation to be made here
is that as the mode number (n) increases, the position of xt moves to the left. Once n is large enough that this point moves
off the membrane (i.e., xt ≤ 0), then the wave for that mode will not propagate. In other words, the higher modes do not
result in propagating waves; these are known as evanescent modes.
The conclusions stated here clearly depend on the problem that was solved. For example, if the membrane widens instead
of tapers, then a completely different response is obtained. Also, in the above example, for each mode there is a possibility
of resonance. This occurs when sin [θ (0) /ε + π/4] = 0, but the analysis of this situation will not be considered.
9 Asymptotic expansion of integrals

9.1 Preamble
The solution of differential equations often leads to formulas involving complicated integrals which cannot be evaluated in
closed form. For example, consider the simple linear differential equation:
′′ ′ ′
y + 2λty = 0; y(0) = 0, y (0) = 1. (654)
′ ′
Letting v = y we obtain v + 2λtv = 0 which gives, on separating variables, integrating, and applying the intial condition
′
y (0) = v(0) = 1, v = exp(−λt2 ). To obtain y we integrate again, now applying the condition y(0) = 0, to get
Z t
y = y(t; λ) = exp(−λs2 )ds (655)
0
We cannot calculate this integral, which represents the solution, because there is no anti-derivative. However, we might want
to know the behaviour of y for large t (with λ fixed), or we may want to determine the value of y for a fixed value of t in the
case that λ is a large parameter. Problems like this are common in applied mathematics, and in this section we introduce a
standard technique to approximate certain kinds of integrals.
9.2 A useful result

R R∞ R∞
exp(−s2 )ds has no anti-derivative but with the specific limits −∞ or 0 we can evaluate the integral in closed form.
Define
Z ∞ Z ∞
I= exp(−x2 )dx = exp(−y 2 )dy (656)
−∞ −∞
so that
Z ∞ Z ∞ Z ∞ Z ∞
2 2 2
I = exp(−x )dx exp(−y )dy = exp(−x2 ) exp(−y 2 )dxdy (657)
−∞ −∞ −∞ −∞
Now transform to polar coordinates via x = r cos θ, y = r sin θ and dxdy = rdrdθ where r is the Jacobian determinant for
the transformation from Cartesians to polars (dxdy = rdrdθ). Thus
Z ∞ Z ∞ Z θ=2π Z r=∞
I2 = exp(−x2 − y 2 )dxdy = exp(−r2 )rdrdθ (658)
−∞ −∞ θ=0 r=0
Now use the substitution r2 = u so that rdr = du/2 and we obtain:

Z Z u=∞ Z
2 1 θ=2π 1 θ=2π
I = exp(−u)du dθ = 1dθ = π (659)
2 θ=0 u=0 2 θ=0
so we get the result:
Z ∞ √
I= exp(−x2 )dx = π. (660)
−∞
132
R∞
Sometimes this is written as −∞ exp(−x2 )dx = Γ(1/2). As the integrand exp(−x2 ) is symmetric about x = 0 it is clear that
Z ∞ √
π
exp(−x2 )dx = . (661)
0 2
A related result is Z r
∞
2 π
exp(−ax )dx = (662)
−∞ a
√
for a > 0. This is easily obtained via the substitution y = ax.
In addition note that:
Z ∞
ta exp(−xt)dt = x−a−1 Γ(a + 1) (663)
0
which can be obtained by repeated integration by parts.
9.3 Expansion of integrands

If the integrand contains a small parameter or if the limits of integration are small one may be able to determine the
asymptotic expansion of the integral by expanding the integrand and then integrating the result term by term.
9.3.1 Example
R1
Expand 0 sint εt dt for small ε.
ε 3 t3 ε 5 t5
sin εt ∼ εt − + + ... (664)
3! 5!
so
Z 1 Z 1 Z 1 Z 1
sin εt ε3 2 ε5 ε3 ε5
dt ∼ ε dt − t dt + t4 dt . . . = ε − + ... (665)
0 t 0 3! 0 5! 0 18 600
The accuracy of this sort of approximation will depend on how many terms one takes and how good an approximation the
Taylor series is for the function over the whole interval of integration. For example, if the interval of integration were [0, 1ε ],
and we use a three term series, the above approach would fail because the series approximation would lose its asymptoticness.
(?)
9.3.2 Splitting the range of integration: leading order estimates

Sometimes it is possible to split the range of integration into a local part, in which some functions may be approximated by
Taylor series, and a global part, which is of lower order. In general we may follow such a procedure, splitting the range of
integration and using different approximations in each range. The major contribution to the integral may come from a local
contribution from a small region or from a global contribution (i.e., from the whole range).
Example 1
Z 1
dx
.
0 (x + ε)1/2
1/2 1/2
The exact answer is 2((1 + ε) − ε ).
On the one hand we would like to expand the integrand for small ε:
1 1 ε
∼ 1/2 − 3/2 + · · ·
(x + ε)1/2 x 2x
However, such an expansion is only okay if ε ≪ x. Thus there are two regions to consider, x = O(1) and x = O(ε). We can
estimate the contribution of an integral as the magnitude of the integrand multiplied by the width of the region. We recall
the inequality:
Z b
f (x) dx ≤ (b − a)max|f (x)|.
a
• If x = O(ε) the integrand is (ε + x)−1/2 = O(ε−1/2 ), the path is O(ε) and the contribution to the integral is therefore
ε
O(ε1/2 ) i.e., ε1/2 = ε1/2 .
133
• If x = O(1) the integrand is O(1), the path is O(1) and the contribution to the integral is therefore O(1).
Thus we expect the global contribution to dominate. From these estimates we can conclude that the leading order term is
a global contribution, for which the integrand can be approximated by x−1/2 and the range of integration is from a small
value outside x = O(ε) which we write as ‘0’ to 1 . Thus
Z 1 Z 1
−1/2
(ε + x) dx ∼ x−1/2 dx = 2
0 ′ 0′
is the leading order contribution. We can verify that this is the correct leading order behaviour by comparing it with the
exact answer.

2 (1 + ε)1/2 − ε1/2 ∼ 2 − 2ε1/2 + ε + · · · .
Example 2
Z π/4
dθ
I= .
0 ε2 + sin2 θ
As before there are two regions, θ = O(1) and θ = O(ε).
• If θ = O(ε) the integrand is O(ε−2 ), the path is O(ε), and the contribution to the integral is therefore O(ε−1 ).
• If θ = O(1) the integrand is O(1), the path is O(1), and the contribution to the integral is therefore O(1).
Thus we expect the local contribution to dominate. To evaluate the local contribution we rescale via θ = εu (so that the
leading contribution comes from u = O(1) as ε → 0). Thus
Z π/4 Z ′
∞′
dθ εdu π
∼ =
0 ε + sin2 θ
2
0 ε2 2
+ε u 2 2ε
d arctan x 1
(recalling the result dx = 1+x2 . The exact answer is
√
1 1 + ε2
√ arctan ( )
ε 1 + ε2 ε
and one can readily verify the leading order estimate by considering the behaviour as ε → 0.
Note An alternative method for evaluating integrals with both local and global contributions is to split the range of
integration into two at some point θ = δ which is large compared with the small region but small compared with the large
region i.e., ε ≪ δ ≪ 1. It may then be possible to find asymptotic approximations to the integral over the two ranges
separately using an appropriate rescaling for the small region. When the two parts of the integral are finally combined the
result should be independent of the artificially introduced δ. It may be useful to tie δ to ε as ε → 0,
√in a way consistent with
ε ≪ δ ≪ 1 i.e., to choose a distinguished limit. In the present example one might choose δ = O( ε). For more details see
Hinch pp 39-40.
9.4 Integration by Parts

In many cases asymptotic expansions of integrals can be obtained by repeated integration by parts. In using this procedure
it is important to check that the “remaining” integral terms are indeed asymptotically small.
9.4.1 Example
To approximate the exponential integral for large x
Z ∞ −t
e
E1 (x) ≡ dt, (666)
x t
we integrate by parts to get
Z ∞ Z ∞ ∞ Z ∞ Z ∞ −t
e−t 1 −t e−t e−x e
dt = − d(e ) = − + e−t d(t−1 ) = − dt
x t x t t x x x x t2
The remainder term is Z ∞ −t
e
R1 (x) = − dt (667)
x t2
134
and we note that because t > x =⇒ 1/x > 1/t in the integral, this term is bounded for large x by
Z
1 ∞ −t e−x
|R1 (x)| ≤ 2 e dt = 2 . (668)
x x x
so the leading order approximation is
e−x e−x
E1 (x) = + O( 2 ) (669)
x x
i.e.
e−x
E1 (x) ∼ . (670)
x
Higher order terms Integrating repeatedly by parts, we find that
e−x e−x e−x

E1 (x) ∼ − 2 + O( 3 ) (671)
x x x
and in general
k
X Z ∞
−x n+1 (n − 1)! k e−t
E1 (x) = e (−1) + (−1) k! dt. (672)
n=1
xn x tk+1
In general the remainder term Z ∞
n e−t
Rn (x) = (−1) n! dt (673)
x tn+1
is bounded for large x by
Z ∞
n! n!e−x
|Rn (x)| ≤ e−t dt = . (674)
xn+1 x xn+1
As a result,
∞
X
−x (n − 1)!
E1 (x) ∼ e (−1)n+1 . (675)
n=1
xn
Notice that although this is an asymptotic representation, it is divergent. Its usefulness, therefore, is not with many terms,
but rather with a small number of terms and x large. Adding more terms eventually makes the approximation worse.
Because E1 (x) ∼ e−x /x then xE1 (x)ex ∼ 1. In the table below we display the function f (x) = xex E1 (x) and its
asymptotic approximation for various values of x. Notice that the best approximation occurs with a different number of
terms for each value of x, and the approximation always deteriorates when many terms are added. However, for a fixed
number of terms, the approximation always improves as x increases.
x xex E1 (x) 2 3 4 5 10
2 0.72265 0.5 1.0 0.25 1.75 −581.37
3 0.78265 0.6667 0.889 0.6667 0.96296 13.138
5 0.85221 0.8 0.88 0.832 0.8704 0.73099
10 0.91563 0.9 0.92 0.914 0.9164 0.91545
20 0.95437 0.95 0.955 0.95425 0.95441 0.95437
100 0.990194 0.99 0.9902 0.990194 0.990194 0.990194
9.4.2 Example
As a second example, consider the incomplete gamma function
Z x
γ(a, x) ≡ e−t ta−1 dt. (676)
0
Recall that the gamma function (which we regard as a known function here)
Z ∞
Γ(a) ≡ e−t ta−1 dt (677)
0
satisfies Γ(a) = (a − 1)Γ(a − 1) and Γ(1) = 1 so that Γ(k) = (k − 1)! if k is a positive integer.
135
X∞
x2 n
x ≪ 1 For small x, we use the fact that e−x = 1 − x + 2! − ... = (−1)n xn! , so that
n=0
Z xX∞ X∞
(−1)n tn+a−1 (−1)n xn+a
γ(a, x) = dt = (678)
0 n=0 n! n=0
n!(n + a)
which converges for all x. In fact the power series converges for all bounded complex numbers. However, its usefulness as an
approximation is limited to small x, since it converges very slowly (it is an alternating series) for large positive x. Of course
the leading order approximation (which is often all one is seeking) is just:
xa
γ(a, x) ∼ (679)
a
x ≫ 1 To find a representation that is useful for large x, we first assume that a < 2 and we define
Z ∞
Γ1 (a, x) ≡ Γ(a) − γ(a, x) = e−t ta−1 dt. (680)
x
Γ(a) is just the value of the gamma function at a and is thus known. If we can evaluate the integral, we can deduce the
value of γ(a, x). We evaluate the integral by parts
Z ∞ Z ∞ ∞ Z ∞ Z ∞
e−t ta−1 dt = − ta−1 de−t = −ta−1 e−t + e−t d(ta−1 ) = xa−1 e−x + (a − 1) e−t ta−2 dt (681)
x x x x x
Z ∞
a−1 −x
=x e + (a − 1) e−t ta−2 dt.
x
The remainder contains:

Z ∞ Z ∞

−t a−2
e t dt ≤ xa−2 e−t dt = xa−2 e−x (682)

x x
as t > x and a < 2 i.e., a − 2 < 0 so our leading order approximation is:
Z ∞
e−t ta−1 dt = xa−1 e−x + O(xa−2 e−x ) (683)
x
or
Z ∞
e−t ta−1 dt ∼ xa−1 e−x =⇒ γ(a, x) ∼ Γ(a) − xa−1 e−x . (684)
x
Higher order terms (and any value of a) Integrating the remainder integral by parts again we see that
Z ∞ Z ∞ ∞ Z ∞ Z ∞
e−t ta−2 dt = − ta−2 de−t = −ta−2 e−t + e−t d(ta−2 ) = xa−2 e−x + (a − 2) e−t ta−3 dt (685)
x x x x x
Z ∞
= xa−2 e−x + (a − 2) e−t ta−3 dt
x
so we have
Z ∞ Z ∞
Γ1 (a, x) = e−t ta−1 dt = xa−1 e−x + (a − 1)xa−2 e−x + (a − 1)(a − 2) e−t ta−3 dt = (686)
x x
Γ(a) −x a−1 Γ(a) −x a−2 Γ(a)
e x + e x + Γ1 (a − 2, x)
Γ(a) Γ(a − 1) Γ(a − 2)
The remainder contains:
Z ∞ Z ∞

|Γ1 (a − 2, x)| = e t −t a−3
dt ≤ xa−4 e−t dt (687)
x x
a−3 −x
=x e
so we have:
136
Γ(a) −x a−1 Γ(a) −x a−2
Γ1 (a, x) = e x + e x + O(xa−3 e−x ) (688)
Γ(a) Γ(a − 1)
and in general we find that:
n
X Γ(a) Γ(a)
Γ1 (a, x) = e−x xa−k + Γ1 (a − n, x). (689)
Γ(a − k + 1) Γ(a − n)
k=1
Now,
Z ∞ Z ∞

|Γ1 (a − n, x)| = −t a−n−1
e t dt ≤ xa−n−1 e−t dt (690)
x x
a−n−1 −x
=x e
for large x, provided n > a − 1. Of course we can always guarantee this by integrating by parts enough times.Thus, we finally
find that
∞
X Γ(a)
Γ1 (a, x) ∼ e−x xa−k (691)
Γ(a − k + 1)
k=1
provided a is not an integer.
9.4.3 Example
As a specific case of the last example, consider the complementary error function for large x,
Z ∞
2 2
erfc(x) ≡ √ e−t dt. (692)
π x
The function occurs frequently in probability theory in connection with the Gaussian distribution. If X is a normally
distributed random variable with zero mean and standard deviation one, then the probability that X > x is P {X > x} =erfc
(x).
If we use the substitution t2 = s, so dt = 2ds
√ then we get
s
Z ∞
1
erfc(x) ≡ √ e−s s−1/2 ds
π s=x2
which is an example of (676) with a = 1/2 < 2. The fact that the lower limit is x2 does not matter as we are looking for
asymptotic behaviour of the integral when its lower limit is large. So the leading order approximation follows from (684)
replacing x by x2 :
Z ∞ 2
1 1 2 1 e−x
erfc(x) ∼ √ e−s s−1/2 ds = √ e−x (x2 ) 2 −1 = √ (693)
π s=x2 π x π
Higher order terms Z

∞
1 −s −1/2 1 1 2
erfc(x) = √ e s ds = √ Γ ,x (694)
π x2 π 2
so that
∞
X
2 1 1
erfc(x) ∼ e−x as x → ∞. (695)
Γ (3/2 − k) x2k−1
k=1
Once again it needs to be emphasized that the computational value of this formula is with one or two terms, not with many
terms, in spite of the notational convenience of an infinite sum.
9.5 Laplace integrals and Laplace’s Method

We define a Laplace integral to be one of the form:
Z b
I(x) = g(t) exp(xh(t))dt, (x → ∞) (696)
a
(where g(t) is continuous, h(t) is sufficiently smooth with a unique maximum (both are known functions)) and we are
interested in the asymptotic behaviour of the integral for x ≫ 1. At the maximum we will have h′ = 0, h′′ < 0 unless
137
the maximum occurs at the endpoint of the interval of interest. In certain cases we can obtain asymptotic approximations
using integration by parts. Note that the method using integration by parts works occasionally but it is inflexible; it can
only produce asymptotic series in integral powers of 1/x and it often fails, for example, when Laplace integrals have large-x
expansions which contain fractional powers of x. (See Bender-Orszag for example).
When integration by parts fails, the workhorse of asymptotic expansions of integrals is Laplace’s method and is based
2
on a very simple observation, namely, if x is large, then an integrand like e−xt is extremely small except near t = 0,
2
the maximum of −t2 . If the integrand contains a factor of e−xt , the main contribution to the integral occurs near the
origin. More generally if the integrand is g(t) exp(xh(t)) then the dominant contribution to the integral comes from some
neighbourhood of the maximum of h(t = t0 ) which can be either an interior or end-point. The fundamental idea is then to
determine what subinterval about t = t0 gives the dominant contribution.
Example The basic idea can be understood by considering the behaviour of the function eMh(t) assuming that h(t) has
a unique global maximum at t = t0 . If we multiply this function by the large positive real number M , the gap between
M h(t) and M h(t0 ) will only increase and it will grow exponentially large when we consider eMh(t) for large M . Significant
contributions to an integral of this function will come only from points in the neighbourhood of t0 and we can use this fact
to simplify the integrand.
Consider the particular example h(t) ≡ sint t which has a global maximum at t = 0. eMh(t) and its approximation based
on expanding h(t) in a series about t = 0 are graphed in fig. 54 for M = 3, 5 illustrating the basic idea. Note that the series
representation about t = 0 is
sin t 1
∼ 1 − t2 .
t 6
sin t
Figure 54: Plots of eM t , M = 3, 5.
In summary, to obtain an approximation for (696) we
1. identify the relevant subinterval about t = t0 where h(t) attains its maximum,
2. approximate the integrand by restricting the region of integration to a narrow region surrounding t = t0 but only
introducing exponentially small corrections to the value of the integral in doing so. Thus replace g(t) by
g(t0 ) and noting that exponentials magnify largeness we replace h(t) by a few terms of its Taylor series about t = t0 ,
3. (if necessary) expand the interval of integration to simplify the computations, ensuring that in doing so the errors
introduced are again always much smaller than the leading order terms in the approximation to the integral. (Why
this is possible is illustrated in example 9.6.2 and figure 56).
The approach here is to develop Laplace’s method by example, and demonstrate how to compute asymptotic approxima-
tions which seem plausible. It may seem foolish to first restrict the region of integration and then expand it! However we
restrict the region of integration to expand the integrand and obtain a series. Having simplified the integrand we sometimes
then re-expand the region of integration if it makes the evaluation of the integral easier.
The ultimate (rigorous) justification that the approximations are asymptotic is via Watson’s lemma.
138
1.0
0.9
0.8
0.7
exp(-xt)
0.6
0.5
x=1
0.4
x=2
0.3
0.2
x=5
0.1
0.0 x=10
0 1 2 3
Figure 55: Illustration of the behaviour of exp(−xt) for increasing values of x. It is clear that when integrating this function,
the main contribution will come from a neighbourhood of t = 0 where −t has a maximum.
9.5.1 The method in more detail

We approximate I(x) by I(x, δ) where
Z t0 +δ
I(x, δ) = g(t) exp(xh(t))dt (697)
t0 −δ
Z a+δ
I(x, δ) = g(t) exp(xh(t))dt (698)
a
Z b
I(x, δ) = g(t) exp(xh(t))dt (699)
b−δ
depending on whether the maximum occurs at an interior point t0 ∈ (a, b) or either of the endpoints. δ is chosen to be
an arbitrary small (fixed) positive number. It is important to remember here that the ultimate limiting process occurs for
x → ∞. There is no limiting process associated with δ; it is a small positive number independent of x. It is also important
to appreciate that the final asymptotic expansion as x → ∞ does not depend on this arbitrary parameter δ
and is identical to the full asymptotic expansion of I(x) as x → ∞. Rather, δ is a device to clarify the procedure. Both of
these results are true because for a < t0 < b,
Z Z
t0 −δ b

g(t) exp(xh(t))dt + g(t) exp(xh(t))dt ≪ I(x) (700)
a t0 +δ
and in fact the error is exponentially small. This follows because for all t on the intervals t ∈ [a, t0 − δ] and t ∈ [t0 + δ, b] ,
exp(xh(t)) is exponentially smaller than exp(xh(t0 )) as x → ∞. To show that I(x) − I(x, δ) is exponentially small as x → ∞,
we use integration by parts. This can be demonstrated in general (exercise) but it is simplest to apply this to each individual
case.
It is helpful to replace I(x) by I(x, δ) because δ > 0 may be chosen so small that it is valid to replace g(t), h(t) by their
Taylor or asymptotic series about t = t0 .
9.6 Example
R∞
Consider the elementary example I(x) = 0
exp(−xt)dt for x ≫ 1. For any δ > 0, in the limit as x → ∞ the following
results hold:
Z ∞
1 1
exp(−xt)dt = − exp(−xt)|∞
0 = (701)
0 x x
Z δ
1 1
exp(−xt)dt = − exp(−xt)|δ0 = (1 − exp(−δx)) (702)
0 x x
139
Z ∞
1 1
exp(−xt)dt = − exp(−xt)|∞
δ = (exp(−δx)) (703)
δ x x
Note that the first two expressions differ only by the exponentially small term exp(−δx) (x → ∞, δ > 0) while in
comparison to them, the last term is itself exponentially small. This simple example contains the essence of the idea behind
Laplace’s method. See also fig. 55. Thus in this case it is clear that I(x) ∼ 1/x. Notice the two basic effects of
increasing x:
1. that the peak of the integrand becomes ever narrower

2. the integrand in the rest of the interval is simultaneously becoming ever smaller.
Similarly it is easy to demonstrate the following result (if g(0) is finite)

Z b
g(0)
I(x) = g(t) exp(−xt)dt ∼ .
0 x
9.6.1 Example
R 10
Let I(x) = 0 (1 + t)−1 exp(−xt)dt so a = 0, b = 10, g(t) = (1 + t)−1 , h(t) = −t. Here h(t) = −t has a maximum at t = 0.
Therefore we replace the integral by:
Z δ
I(x, δ) = (1 + t)−1 exp(−xt)dt (704)
0
for any δ > 0 at the cost of introducing errors which are exponentially small as x → ∞. Next we choose δ so small that we
can replace (1 + t)−1 in the integrand by 1, i.e., the first term of its Taylor series about t = 0. This replacement makes the
integral easy to evaluate. In this particular example, there is now no necessity to look for an expansion of h(t) = −t because
it is already such a simple function. Thus
Z 10 Z δ
−1 1 − exp(−δx)
(1 + t) exp(−xt)dt ∼ exp(−xt)dt = (x → ∞) (705)
0 0 x
Since exp(−δx) ≪ 1 as x → ∞ for any (fixed) δ > 0, we obtain
Z 10
I(x) = (1 + t)−1 exp(−xt)dt ∼ 1/x, (x → ∞). (706)
0
Note that the final result does not depend on the arbitrary parameter δ; it appears only in an exponentially small correction
term (exp(−δx) is exponentially small as x → ∞). Note also that in this simple example, it was not necessary to expand the
interval of integration in I(x, δ). The problem was simple enough to permit a closed form solution while integrating over the
interval [0, δ].
Higher order terms To obtain higher order terms, we can replace (1 + t)−1 by its full Taylor expansion so that (1 + t)−1 =
P
(−t)n which converges for |t| < 1. Thus
Z 10 Z δ ∞ Z
X δ
(1 + t)−1 exp(−xt)dt ∼ (1 + t)−1 exp(−xt)dt = (−t)n exp(−xt)dt (707)
0 0 n=0 0
Now the easiest way to evaluate this integral is to expand the region of integration
R∞ from [0, δ] to [0, ∞). In doing so
we introduce only exponentially small terms as x → ∞ because the integral δ is exponentially small with respect to
R∞ n n −n−1
0 (−t) exp(−xt)dt = (−1) n!x as x → ∞.
We verify this using integration by parts:
Z ∞
(−t)n exp(−xt)dt ∼ (−δ)n exp(−δx)/x, (x → ∞) (708)
δ
R∞
which is indeed exponentially smaller than 0 (−t)n exp(−xt)dt as x → ∞. Assembling the results we obtain:
Z 10 ∞
X
−1
(1 + t) exp(−xt)dt ∼ (−1)n n!x−n−1 , (x → ∞). (709)
0 n=0
140
9.6.2 Example
Z π
Approximate I(λ) = exp(λ sin t)dt for large λ.
0
Note that sin t has a maximum at t = π/2 so we express it in a Taylor series about t = π/2. If you are used to writing
Taylor series at the ‘new point’ x + h in terms of quantities evaluated at the ‘old’ point x:
df
f (x0 + h) = f (x0 ) + h (x0 ) + O(h2 )
dx
then let x0 = π/2 and x0 + h = t =⇒ h = t − π/2 to get
f (t) = f (π/2) + (t − π/2)f ′ (π/2) + . . .
Hence
π π π 1 π π π
sin(t) ∼ sin + (t − ) cos − (t − )2 sin . . . = 1 − (t − )2 . . . (710)
2 2 2 2 2 2 2
Thus we write:
Z π Z π/2+δ Z π/2+δ Z ∞
1 π 1 π
exp(λ sin t)dt ∼ exp(λ sin t)dt ∼ exp(λ(1 − (t − )2 )dt ∼ exp(λ(1 − (t − )2 )dt (711)
0 π/2−δ π/2−δ 2 2 −∞ 2 2
√
2π exp(λ)
= √ (712)
λ
R∞ √
Here we have used the result −∞ exp(−s2 )ds = π.
As we gain more experience, we may relax the technique by not going through the formal procedure of introducing δ > 0.
Figure 56: Illustration of the exact integrand (λ = 10) and the approximate integrand as applied to example 9.6.2. It is
important to note that while the approximate integrand exp(λ(1 − 12 (t − π2 )2 ) is approximately equal to the exact integrand
exp (λ sin t) about t = π/2, away from this point the two integrands are quite different. In fact the approximate integrand is
exponentially small everywhere except in a neighbourhood of t = π/2. This feature of Laplace’s method is what allows
one to change the interval of integration toR(−∞, ∞) if this makes evaluation of the integralReasier. Note that this is clearly
π ∞
not the case for the exact integrand i.e., 0 exp (λ sin t) is NOT approximately equal to −∞ exp (λ sin t) because of the
periodicity of sin t.
9.6.3 Example
Consider (for large fixed x) the integral
Z ∞
ψ(x) = exp(−x cosh θ)dθ. (713)
−∞
The integrand is a “bell-shaped” curve of height e−x whose width decreases as x increases. It has a maximum at θ = 0.
Since the integrand is inconsequentially small for large θ and large x we replace cosh θ by its Taylor series about θ = 0
141
i.e.,cosh θ ∼ 1 + θ2 /2 which for small θ is accurate, but for large θ, where it does not matter, is in error. With this
substitution we calculate that
Z ∞ Z ∞ 1/2
2 −x −xθ 2 /2 −x 2π
ψ(x) ∼ exp(−x(1 + θ /2))dθ = e e dθ = e (714)
−∞ −∞ x
9.6.4 Example
Another example is to approximate the gamma function
Z ∞
Γ(x + 1) = e−t tx dt (715)
0
for large fixed x. In its current form, it is not obvious how to apply Laplace’s method. However, if we note that
tx = exp(ln tx ) = exp(x ln t) (716)

we can write the integral as:
Z ∞
Γ(x + 1) = ex ln t−t dt (717)
0
and observe that the function x ln t − t has a maximum at t = x, then the change of variables t = xs is suggested.
With this change of variables
Z ∞
Γ(x + 1) = xx+1 e−x(s−ln s) ds (718)
0
is now in a form to which Laplace’s method is applicable. The function −s + ln s has a maximum at s = 1 and is concave
downward away from s = 1. Therefore, the integrand is ”bell-shaped” (although not symmetric) and for large x is exceedingly
narrow. Thus, we replace the function s − ln s by its local approximation about s = 1
(1 − s)2
s − ln s ∼ 1 + (719)
2
and obtain
Z ∞
2
Γ(x + 1) ∼ xx+1 e−x e−x(1−s) /2
ds. (720)
0
We still cannot evaluate this integral, but for large x it can be evaluated by changing the interval of integration from [0, ∞)
to (−∞, ∞). (This introduces
√ an error which is exponentially small because the maximum occurs at s = 1). Thus (using the
substitution s − 1 = 2τ ) the integral is approximated by
Z ∞
2
Γ(x + 1) ∼ xx+1 e−x e−x(1−s) /2
ds (721)
−∞
1/2
x+1 −x 2π
=x e .
x
This approximation of Γ (x + 1) is known as Stirling’s formula.
9.7 Discussion
The idea of Laplace’s method should by now be clear. We seek an asymptotic expansion for large x for the integral
Z b
f (x) = exh(t) g(t)dt
a
where h(t) has a local maximum at t = t0 in the interval [a, b]. We can approximate h(t) by
(t − t0 )n
h(t) = h(t0 ) + h(n) (t0 ) + ···
n!
where h(n) (t0 ), the nth derivative of h with respect to t at t0 , is negative, and n is even.
A first order approximation to f (x) is found by replacing h(t) by
142
(t − t0 )n
h(t0 ) + h(n) (t0 ) ,
n!
g(t) by g(t0 ), and (if necessary) expanding the limits of integration to the infinite interval (−∞, ∞) . (This introduces an
error which is exponentially small). From this approximation we obtain
Z ∞ xhn (t0 )(t−t0 )n

f (x) ∼ exh(t0 ) e n! g(t0 )dt
−∞
!1/n
n! 1
= 2g(t0 ) Γ + 1 exh(t0 ) .
x h(n) (t0 ) n
Higher order approximations can be obtained in principle by retaining more terms in the expansions for g(t) and h(t). This
is not as simple as one might expect (see e. g., Bender and Orszag which contains many advanced examples). In particular,
one must be careful with the expansion of h(t) if one goes beyond the second order.
Watson’s lemma The statement that this technique (Laplace’s method) is valid is known as WATSON’S LEMMA which
we now state.
Consider (696) with h(t) = −t, a = 0 so that
Z b
I(x) = g(t) exp(−xt)dt, b>0 (722)
0
so h(t) = −t has a maximum at t = 0. Suppose g(t) has an asymptotic series expansion for small t (this is usually but not
necessarily a Taylor series)
∞
X
g(t) ∼ tα an tβn , (t → 0+ ) (723)
n=0
This includes the case where g(0) is finite via α = 0, a0 = g(0) but also allows for cases where g(t) does not have a Taylor
series about t = 0. Note that we must require that α > −1 and β > 0 for the integral to converge at t = 0. In addition, if
b = +∞, it is necessary that g(t) ≪ exp(ct), (t → +∞), for some positive constant c, in order for the integral to converge.
Watson’s lemma then states that under these conditions:
X∞
an Γ(α + βn + 1)
I(x) ∼ , (x → +∞) (724)
n=0
xα+βn+1
The leading order result (the first term in the series) is what we would obtain if we used Laplace’s method as described in
the preceding pages supposing that g(0) exists finitely. In this case in (724) we have α = 0 and a0 = g(0). Then (724) yields
the leading order approximation:
a0 Γ(1) g(0)
I(x) ∼ = .
x x
N.B. Watson’s lemma essentially allows us to substitute in a series expansion for g(t), change the limits of integration to ∞
and integrate term by term using the result (707). A proof is given for example in Hinch p. 27.
Watson’s lemma only applies to Laplace integrals with the particular form h(t) = −t. For more general h(t) there are two
possible approaches. If h(t) is sufficiently simple, it may be useful to make a change of variable by susbtituting
s = −h(t) (725)
R −h(b)
into (696) and to rewrite the integral in the form −h(a) F (s) exp(−xs)ds where F (s) = − gf′(t)
(t) . Watson’s lemma applies to the
transformed integral. Sometimes such a substitution is unwieldy because the inverse function t = g −1 (−s) is a complicated
multi-valued function, In this case, it may be simpler to use a more direct method than Watson’s lemma for finding the first
few terms in the asymptotic series for I(x).
Example Expand (for large x)

Z 5
exp(−xt)
I(x) = dt. (726)
0 (1 + t2 )
For large x we replace g(t) = (1 + t2 )−1 by its Taylor series about t = 0
(1 + t2 )−1 = 1 − t2 + t4 − t6 + . . . (727)
143
Watson’s lemma allows us to substitute the Taylor series into the integral, interchange orders of integration and summation,
and replace the upper limit of integration by ∞. This gives:
1 2! 4! 6!
I(x) ∼ − 3 + 7 − 7 + . . . (x → +∞). (728)
x x x x
Example Approximate for large x:

Z π/2
I(x) = exp (−x sin2 t) dt. (729)
0
If we want to rigorously find an infinite series expansion, Watson’s lemma cannot be applied directly as h(t) = − sin2 t.
However if we use the substitution s = sin2 t, the integral becomes:
Z
1 1 1
I(x) = exp (−xs) ds, (730)
2 0 s(1 − s)
we can directly apply Watson’s lemma since:
1 1 1 3
= √ + s + s3/2 + · · · (731)
s(1 − s) s 2 8
and we can integrate term by term. The leading order approximation is
√
1 π
I(x) ∼ √ . (732)
2 x
Alternatively, we can apply Laplace’s method directly (without the rigorous underpinning of Watson’s lemma). We note
that the maximum of − sin2 t on [0, π/2] occurs when t = 0. We thus replace sin2 t with its Taylor series about t = 0:
1
sin2 t ∼ t2 − t4 + . . .
3
Retaining the leading order term we find that
Z ∞ √
−xt2 1 π
I(x) ∼ e dt = √
0 2 x
on using (662).
9.8 Method of Steepest Descents (Saddle point method)

The asymptotic evaluation of contour integrals centres around one question and that is how best to deform the contour
of integration to take advantage of Watson’s Lemma. Let us consider integrals of the form:
Z
I(λ) = g(z) exp (λh(z)) dz (733)
C
where λ ≫ 1 is a large positive parameter, g(z), h(z) are known analytic functions of z = x + iy in some domain in the
complex plane which contains C, the integration path. In practice, we may have to choose the region to avoid singularities
or branch points of the functions g(z), h(z).
Summary of the method (A good review of complex integral analysis is to be found in Greenberg’s Advanced Engineering
Mathematics). By Cauchy’s theorem, the path of integration may be replaced by another path C ′ , having the same endpoints
as the original path, provided that C can be continuously deformed into C ′ without leaving the domain of analyticity of
the integrand. If the endpoints of C and C ′ are P1 , P2 , then we can choose arbitrarily a path joining these endpoints and
forming
H a closed loop L in the complex R plane.
R If the integrand is analytic in this closed loop, then Cauchy’s theorem tells us
that L = 0 and we can deduce that C = C ′ .
We will thus attempt to deform our path of integration joining P1 and P2 in such a way that, for large λ, the major part
of the contribution to the integral arises from a small portion of the new path of integration where we can apply Laplace’s
method to obtain an asymptotic approximation for the integral. More specifically, the idea of the method is to use
the analyticity of the integrand to justify deforming the path C to a new path C ′ on which h(z) has a constant
imaginary part. Once this has been done, I(λ) may be evaluated asymptotically using Laplace’s method. To
see why observe that on the contour C ′ we may write h(z) = φ(z) + iψ where ψ is a real constant and φ(z) is a real valued
function of z = x + iy. Thus the integral takes the form:
Z
I(λ) = exp (iλψ) g(z) exp (λφ(z)) dz.
C′
Although z is complex, the integral can be treated using Laplace’s method because φ(z) = ℜh(z) is real. There is no rapid
oscillation (and possible cancellation) effect because ψ =constant along the path chosen.
144
C
z=b
z=a
C'
Figure 57: Cauchy’s theorem.
The method in more detail

Note first that if we write
h(z) = φ(z) + iψ, φ = ℜ(h(z)), ψ = ℑ(h(z)) (734)
then
e(λh(z)) = eλφ(z) eiλψ(z) = eλφ(z) (cos λψ(z) + i sin λψ(z)). (735)
So e(iλψ(z)) has an oscillatory component which can cause difficulties to estimates of the integral due to cancellations etc. The
main contribution from the integral in absolute value comes from eλφ(z) so the details of the “topography” of the function
of two variables φ(x, y) play an important role in what follows.
Consider now a limit inequality argument for I(λ). If the contour C, which may be finite or infinite, joins the points
z = a and z = b, then we have from (733) and (734)
Z sb Z sb
λh(z)
|I(λ)| ≤ |g(z)e |ds ≤ |g|eλφ ds (736)
sa sa
where we have parameterized the integral using arclength (ds q = |dz|), sa and sb correspond to the end points z = a and
z = b of the path and we have used that fact that |e (iλψ)
| = cos2 (λψ) + sin2 (λψ) = 1. If we assume for this preliminary
R R
discussion that C g(z)dz is absolutely convergent i.e., C |g(z)|dz converges, then for large λ the integral in (736) is O(eλφ ),
1
except for multiplicative algebraic terms like λ− 2 , λ−1 etc. If the contour is of finite length L, say, then in place of (736) we
have
|I(λ)| < L maxL |g(z)|eλz . (737)
In (736) and (737) the most important contribution to the asymptotic approximation for |I(λ)| as λ → ∞ must come from
the neighbourhood of the point of maximum φ. We now exploit the fact that the path of integration can be deformed, by
using Cauchy’s theorem.
Recall that Cauchy’s theorem states that I
f (z)dz = 0 (738)
L
for all closed loops L provided f (z) is analytic inside L. In particular if we have a path of integration
R R between two points
z = a and z = b, and we form a second (non-intersecting) path C ′ joining z = a and z = b, then C = C ′ if f (z) is analytic
inside the closed loop formed by C and C ′ (see fig. 57). In fact, the line integral of an analytic function around any simple
closed curve is equal to the line integral of the same function around any other simple closed curve into which the first can
continuously be deformed without passing through a point where f (z) is not analytic.
145
Comment We use this fact to deform C. If g has an isolated pole singularity, for example, we can still deform the path
into another which involves crossing such a singularity, if we use the theory of residues appropriately. Branch points and
cuts are more complicated to deal with.
We now deform C so that it not only passes through the point z = z0 , say, where φ = ℜ(h(z)) has its maximum value,
but also along it φ drops off on either side of the maximum as rapidly as possible. In this way the largest value of eλφ will be
concentrated in a small section of the contour. This specific path through the maximum of eλφ will be the path of steepest
descents. In the case of (737) the length of the path also varies in this contour deforming exercise. When a path is close to
the optimal one in the steepest descents sense, a very small variation in the path, and hence its length L, can drastically
change the variation of eλφ in the neighbourhood of its maximum when λ is large. Thus, if we are interested only in |I(λ)|
as λ → ∞ we would choose a path which made φ and and hence eλφ behave in the above manner, and then use Laplace’s
method. However, since we are interested in I(λ) and not just its modulus as λ → ∞, we must be more specific, keeping in
mind the above discussion.
If we take an arbitrary path through z0 , the point giving the maximum φ, the imaginary part of ψ of h(z) gives an
oscillatory contribution of eiλψ to the integrand. These oscillations are increasingly more rapid the larger λ becomes. This
jeopardises the whole procedure unless a path can be chosen which reduces the effect of the oscillations on the integrand in
the vicinity of z0 , so that we can use the above idea. An appropriate path to choose to overcome this problem is one in which
ψ = ℑ(h(z)) =constant, in the vicinity of z0 , so that there are no oscillations from eiλψ near the point of maximum eλφ .
Returning to (733) we recall:
z = x + iy; h(z) = φ(x, y) + iψ(x, y). (739)
To search for a turning point of the function of two variables φ, we require that ∇φ = 0. Since φ and ψ are the real and
imaginary parts of h(z), an analytic function of z, they satisfy the Cauchy-Riemann equations so
∂φ ∂ψ ∂φ ∂ψ
= , =− =⇒ ∇2 φ = 0; ∇2 ψ = 0. (740)
∂x ∂y ∂y ∂x
If z0 = x0 + iy0 is a point at which ∇φ = 0, then z0 is also a solution of ∇ψ = 0 and so:
h′ (z) = φx + iψx = φx − iφy = 0 when z = z0 = x0 + iy0 . (741)
But h(z) is analytic so φ and ψ both satisfy Laplace’s equation i.e.,
∇2 φ = φxx + φyy = 0; ∇2 ψ = ψxx + ψyy = 0. (742)
Hence φ, ψ cannot have a maximum (or a minimum) in the domain of analyticity of h(z) as this would require φxx < 0 and
φyy < 0 with a similar comment for ψ. The turning point z0 is thus a saddle point of both φ, ψ and |h(z)| (see figure 58).
Thus we say also that z0 is a saddle point of h(z). Here we will restrict attention to saddle points of order one where
h′ (z0 ) = 0, h′′ (z0 ) 6= 0. (743)
The technique developed here can be extended for saddle points of order n. If there are no saddle points h(z) must be a
simple linear function and the method concerns itself with the contributions from the endpoints and choosing a suitable path
through these endpoints. One can compare this to finding the maximum or minimum of a function of a single real variable
h(t). In the case where there are no points where h′ (t) = 0, one has to examine the endpoints. For example, h(t) = −t on
the interval [0, b) has a maximum at t = 0 and a minimum at t = b. Use of Laplace’s method for evaluating the definite
Rb
integral 0 would involve focussing on the region in the neighbourhood of t = 0.
If we consider the surface given by φ = φ(x, y) in x, y, φ space, a typical saddle point situation is illustrated in figure
58 and 59 where there is both a surface plot and a contour plot in the xy plane. (Recall the condition for a saddle point:
φxx φyy − φ2xy < 0.) Here the point S ′ in the surface is the saddle point corresponding to the point S at z = z0 in the xy
plane. The contour diagram of fig. 59 thus represents the projections onto the z−plane of the intersections of the planes
φ =constant and the three dimensional surface surface φ = φ(x, y) in the φ, x, y space. The contours, as usual, form a 2D
representation of the surface in that they are lines joining points of equal height on the surface. Primes on the points in the
surface relief correspond to the unprimed points in the z−plane. From (741) the tangent plane at the saddle point S ′ on the
surface is a constant φ plane parallel to the z−plane.
We return to the question of the optimal path into which we must deform the path C. To avoid the oscillations of eiλψ ,
we will choose a path along which ψ = ℑh(z) =constant. From the Cauchy-Riemann equations (740), since
∂φ ∂ψ ∂φ ∂ψ
∇φ · ∇ψ = + = ψx ψy − ψx ψy = 0, (744)
∂x ∂x ∂y ∂y
the normals to the level curves of φ and ψ are orthogonal. It follows that the (level) curves of constant φ and ψ must
themselves be orthogonal to one another. From the definition of the directional derivative it follows that
dψ
= n̂ · ∇ψ = 0 (745)
dn
146
Q'
P'
S'
N'
M'
Figure 58: 3D plot of the surface representing a typical saddle point in the complex plane.
Q
S
P
M
Figure 59: Contours of a typical saddle point in the complex plane.
147
when n̂ is in the direction of ∇φ i.e., ∇φ is in the direction along which ψ =constant and vice versa. (∇φ and ∇ψ are vectors
pointing in the direction of greatest change of φ and ψ respectively).
Thus the curves along which φ changes most rapidly, that is the direction of ∇φ, are the curves ψ =constant. If we
choose the constant ψ−curve which passes through the point z0 , where φ has its saddle-point, then this is in keeping with
the necessary conditions that (i) along the optimal path φ has as rapid a variation as possible near its relative maximum
and (ii) there are no oscillation contributions from eiλψ . Looking in detail at figures 58 and 59 we see from the sketch of the
contour lines of constant φ in the z−plane that there are two curves ψ =constant which pass through the saddle point at
z0 and along which φ changes as rapidly as possible. However P SQ corresponds to the steepest ascent path (P ′ S ′ Q′ on the
surface) since the value of φ on this path satisfies φ(x, y) > φ(x0 ) except at z0 . If we instead consider the curve of constant
ψ passing through M SN in the z−plane, we see that every point on it satisfies φ(x, y) < φ(x0 , y0 ), except at z0 . Thus the
curve ψ =constant through M SN is the line of steepest descents and the one into which C should be deformed.
It might be argued that a path which starts and finishes on either side of the ridge of the saddle-point but which passes
through a higher point is better than the one chosen. However, in this situation we do not stay on a curve of constant ψ
and we have an oscillation contribution from eiλψ and the argument that the major contribution to the integral as λ → ∞
comes from the region of maximum φ is no longer valid.
The correct path of the two ψ =constant lines through the saddle point is obtained simply by considering the values of φ
along both and choosing that on which φ(x, y) > φ(x0 , y0 ) except at the saddle point. We also note that if there are several
saddles through which the contour must pass, as it progresses across various ridges from one fixed end point to the other,
then the highest of these will dominate the integral up to exponentially smaller terms as z → ∞.
N.B. In practice when solving difficult problems one does not usually have to find the actual path of steepest descent from
the saddle point: one is usually only interested in a leading order or one or two term approximation and any path descending
from the saddle on both sides will produce an integral which converges locally to the correct answer. (Our first example
illustrates this point: we first find a leading order approximation without finding the path of steepest descent. Then, using
the path of steepest descent we find the full asymptotic expansion).
Note also that if a particular h(z) has no saddle points (e.g., h(z) = z has no saddle points as h′ (z) = 1 everywhere), the
resulting surface is planar and we must look for leading order contributions to the integral by deforming the contour along a
curve of constant ψ passing through the endpoints. The dominant contribution will come from a neighbourhood of one of the
endpoints. (Compare the fact that for a linear function of a single variable (f (x) = mx + C), the maximum (and minimum)
occurs at an endpoint and f ′ (x) 6= 0 ). Thus, unlike the case of a saddle point where f ′ (z) = 0, the series expansion that
one takes near the maximum will have a term involving the first derivative.
Example
Consider the integral Z
1 √ dz
I (λ, a) = exp λ za − z
2πi C z
where a is real, positive, and the contour C is any vertical line of the form z = xc + iy, xc > 0, fixed, −∞ < y < ∞. It is
sufficient to evaluate this integral with a √
= 1 since for any other value of a, I (λ, a) = I (λ/a, 1). We drop the reference to a
and set I(λ) ≡ I (λ, 1). Note that ℜ(z − z) < 0 at the end points in order for the integral to converge. This ensures that
integrating along the arc at ∞ yields no contribution. The value of this integral is clearly independent of the value of xc
provided xc > 0. To know how best to deform the path C, we must first understand the structure of the function
√
h(z) = z − z.
We need to find its saddle point and a line of constant ψ through this saddle point.
A leading order approximation Note that h(z) has a branch point4 at z = 0 and we take the branch cut5 to be along
the negative real z axis. At z0 = 14 , the function h(z) has zero derivative, h′ (z0 ) = 0 because
dh 1 √ 1 1
= 1 − √ = 0 =⇒ z = =⇒ z = z0 = .
dz 2 z 2 4
The point z0 is a SADDLE POINT of |h(z)| (analytic functions cannot have local maxima or minima).
The idea of
√ the method of steepest descents is to deform the path of integration into one for which the real part of the
function z − z decreases as rapidly as possible. The hope is that, for λ large, the integral is well approximated by the
integrand near the saddle point and the rest can be ignored with impunity. In other words, we want to choose a path along
which Laplace’s method can be applied.
4 If different values of a function f (z) are obtained by successively encircling some point z in the complex plane, then the point z is called a
0 0
branch point.
5 A branch cut is a barrier in the complex plane introduced to make a multi-valued complex function single valued.
148
It is easiest to look at the behaviour of h near the saddle pointqand to choose a path on which ℜ(h) decreases as fast as
possible initially. At the saddle point z0 = 14 , we have h( 14 ) = 41 − 14 = − 41 so near the saddle point we can express h(z) in
a Taylor series about this point, noting that:
d2 h 1 1 3 1
2
= − (− )z − 2 = 3 .
dz 2 2 4z 2
Thus 2
1 1 ′ 1 1 1
h(z) = h( ) + (z − )h ( ) + . . . = − + z − + ...
4 4 4 4 4
3
as ( 41 ) 2 = 18 .
Note that because
√ z0 = 41 is a saddle point, h′ (z0 ) = 0 so the term in the Taylor series proportional to z − 14 is identically
zero. (Though z does not have a Taylor series about z = 0 it does have one about z = 41 ).
Clearly, to make ℜ(h) decrease we can choose a path of integration xc + iy along which xc = 14 i.e., we choose a path
along which z = 41 + iy, dz = idy for −δ < y < δ. If we make this change of variables,6 and thus approximate h(z) by
1 1 1
h(z) ∼ − + (z − )2 ∼ − − y 2 , g(z) ∼ 4,
4 4 4
we find that
Z
1 −λ/4 δ exp −λy 2 dy 1 1
I(λ) ∼ e 1 and using the fact that ( + iy) ∼ + . . .
2π −δ 4 + iy 4 4
Z δ Z ∞
2 2
∼ e−λ/4 exp −λy 2 dy ∼ e−λ/4 exp −λy 2 dy
π −δ π −∞
1/2
1
=2 e−λ/4 .
πλ
Here we have used Laplace’s method and the fact that on the integration path
1 1
( + iy)−1 ∼ .
4 4
The approximation we have just made is the first term of an expansion using Watson’s lemma.
A more accurate approach: more terms If one requires more information, an estimate of the error of this approx-
imation, or a rigorous statement of its validity, it is necessary to find more terms. To do so we need a more accurate
representation of the steepest descent path.
How do we find the steepest descent path? From calculus, we know that for a real valued function, the steepest paths
are along gradient directions. For a complex analytic function h(z) = φ + iψ, ∇φ·∇ψ = 0 because of the Cauchy-Riemann
conditions. In other √words, curves with ℑ(h) constant are steepest paths for ℜ(h) and vice versa. Thus we have
h(z) = φ + iψ = z − z and we wish to find paths along which ψ = 0. Using polar coordinates z = reiθ we have:
√ 1 iθ 1 θ 1 θ
h(z) = z − z = reiθ − r 2 e 2 = r cos θ − r 2 cos + i(r sin θ − r 2 sin ) = φ + iψ.
2 2
1
The curve of constant ψ passing through the saddle point z = z0 = 4 satisfies h(z0 ) = − 41 so ψ = 0 along this curve. To
solve this equation we write it as:
θ 1
1

= 0 =⇒ sin θ/2 2r cos θ/2 − r 2 = 0.
2r sin θ/2 cos θ/2 − r 2 sin
2
θ
Thus one solution to ψ = 0 is given by sin 2 = 0. This corresponds to θ = √
0, 2π, 4π . . . i.e., the positive x axis. This is not a
path of steepest descent of h(z) through z0 = 1/4 as, in fact, h(z) = x − x is increasing as x moves away from 1/4. The
other solution of ψ = 0 is
1 1
r2 = .
2 cos θ2
Thus
1 1 1
r= θ
= .
4 cos2 2
2 (1 + cos θ)
6 As 1
this is a definite integral, there is no necessity to parameterise the path of integration though one could also do this via: x = 4
, y = τ,
−δ ≤ τ ≤ δ, z = x + iy = 14 + iτ , dz = idτ .
149
2
C Re(f) = -1/4
Im(z)
0
Im(f) = 0
-2
-2 0 2
Re(z)
Figure 60: Paths of steepest descent and ascent for h(z) = z − z 1/2 .
p
Reverting to cartesian coordinates we have θ = arctan xy , r = x2 + y 2 and we find that the path of steepest descent,
expressed in cartesian coordinates, is a parabola:
1 1 1
r= =⇒ x = − y 2 , (746)
2 (1 + cos θ) 4
and we verify below that φ does decrease if we follow this path i.e., the path does indeed descend from the saddle point.
Though we do not use it here, we can represent the steepest ascent curve φ = ℜ(h) = − 41 in polar coordinates z = reiθ ,
and we find
1 1
r=
4 1 ± sin θ
are the two branches (shown in √ figure 60).
For the function h(z) = z − z the curves φ = ℜ(h) = − 41 and ψ = ℑ(h) = 0 are shown in Figure 60.
The parabola (746) is the exact path into which we wish to deform the path of integration C. Notice that deformation is
possible since the arcs at infinity do not contribute to the integral. (As the original integration path is infinite, the integrand
must go to zero at the endpoints xc ± i∞ in order for the integral to converge. In fact ℜ(h) is large and negative at every
point in the region of deformation - use the idea behind Jordan’s Lemma.)7
Along the chosen integration path x = 1/4 − y 2 , we note that on substituting x = 41 − y 2 into z = x + iy we obtain:
1 1
z = x + iy = − y 2 + iy = ( + iy)2 .
4 2
1
To find a parametric representation of the path of integration we simply choose y = τ , x = 4 − τ 2 and so
2
1
z = x + iy = + iτ
2
and dz = (i − 2τ )dτ = 2i( 21 + iτ )dτ where τ is real and −∞ < τ < ∞. This corresponds to the curve ℑ(h) = ψ = 0. Along
the path of integration we note that
2
√ 1 1
h(z) = z − z = + iτ − + iτ = −1/4 − τ 2 .
2 2
7 Let
R
I(R) ≡ Γ(R) eiaz g(z)dz where g(z) is single-valued, a is real, and the path of integration coincides with the real axis except for indentations
at the singularities of g(z). We take a > 0; the case a < 0 is similar. In the treatment of such integrals one is frequently concerned with
Z π
I= ieiaR cos θ−aR sin θ g(Reiθ )Reiθ dθ.
0
In many such integrals Rg(Reiθ ) does not vanish as R → ∞. We note, however, that for |g(Reiθ )| ≤ G(R),
Z π Z π
πG(R)
|I| ≤ e−aR sin θ RG(R) dθ ≤ 2 RG(R)e−2aRθ/π dθ = (1 − e−aR )
0 0 a
and I → 0 as R → ∞ provided only that G(R) → 0 as R → 0. This result is Jordan’s lemma; it is broadly useful, not only in the context given
above, but also in connection with integrals for which g(z) is multivalued and for which isolated pieces of the semicircle are under consideration.
150
This clearly corresponds to a path along which |h(z)| decreases as we move away from the saddle point at which point
h = −1/4, τ = 0. It is also simple enough for the purposes of integration. In addition, along the path of integration we have:
1 1
g(z) = = 1 .
z ( 2 + iτ )2
Every term in this change of variables is well√motivated by geometrical reasoning. We want a curve parameterized by τ,
passing through z = 1/4 when τ = 0. Thus, z√ − z = −1/4 when τ = 0. We want a steepest descent path so we require τ to
be real since then the imaginary part of z − z is constant. We pick −τ 2 rather than +τ 2 dependence to be sure the real
part becomes negative (and hence we get a bell-shaped integrand along the path) and we take quadratic dependence on τ
since the saddle point has locally quadratic behaviour. In terms of the new variable τ note that with g(z) = 1/z, we have
dz/z = ( 12idτ and the integral simplifies to:
2 +iτ )
Z
2 −λ/4 ∞ exp(−λτ 2 )dτ
φ(λ) = e
π −∞ 1 + 2iτ
and, of course, −τ 2 has a maximum at τ = 0.
Although this expression is exact, we cannot evaluate it exactly. Proceeding directly we could use an integral substitution
to deal with the exp (−λτ 2 ) in the integrand. It is also possible to extend Watson’s lemma to deal with integrands containing
−1
exp (τ n ) where n is any positive integer. Using the extended form of the lemma, we expand the function (1 + 2iτ ) in a
power series for small τ, keeping as many terms as we desire, and then evaluate the integrals term by term. This is guaranteed
(rigorously!) by Watson’s lemma to yield the asymptotic result
1/2
1 −λ/4 2 12 120 1
φ(λ) ∼ 2 e 1− + 2 − 3 +O .
πλ λ λ λ λ4
We can use the geometric series
∞
X
−1 n
(1 + 2iτ ) = (−2iτ )
n=0
to write the full asymptotic approximation as

∞
2 −λ/4 X Γ(n + 1/2)(−4)n
φ(λ) ∼ e .
π n=0
λn+1/2
9.8.1 Example Hankel’s integral

Consider
Z
1 1
= ez1 z1−γ dz1
Γ(γ) 2πi C
where the contour C traverses from z1 = −ia − ∞ to z1 = ia − ∞ around a branch cut on the ℜ(z1 ) < 0 axis, a > 0, (Figure
61).
To approximate this integral for γ ≫ 1 we first make a change of variables z1 = γz which turns a moveable saddle into a
fixed one. Thus
Z
1 1 1−γ
= γ eγ(z−ln z) dz.
Γ(γ) 2πi C
d 1
The function f (z) = z − ln z has a saddle point when dz (z − ln z) = 1 − z = 0 i.e., when z = 1, f (1) = 1. f (z) can be
represented as a power series near z = 1 by
X∞
(1 − z)n
f (z) = 1 + .
n=2
n
The steepest descent path through the saddle point z = 1 satisfies ℑ(z − ln z) = 0. Writing z − ln z = φ + iψ we have, on
writing z = x + iy
p y p y
z − ln z = x − ln ( x2 + y 2 ) + i(y − arctan ) =⇒ φ = x − ln x2 + y 2 , ψ = y − arctan ,
x x
on recalling that ln z = ln r + iθ in polar coordinates. Hence, the path of steepest descent is given by:
y
ψ = constant = y − arctan .
x
151
5
C' π
Im(z)
0
−π
-5
-5 0 5
Re(z)
Figure 61: Path of steepest descent for Hankel’s integral.
We require the path to pass through the saddle at z = 1 i.e., x = 1, y = 0 so the constant is zero and the steepest descent
path is:
y
y − arctan = 0 =⇒ x = y cot y
x
which is C ′ in figure 61.
Leading order approximation At leading order we merely need to pass locally along the steepest descent path through
the saddle point. Thus it is sufficient to consider the path C ′′ parallel to the y−axis from z = 1 − iδ to z = 1 + iδ. Near the
saddle we use a Taylor series about z = 1 (cf. (2.2.1)
(z − 1)2 (z − 1)3
z − ln z ∼ 1 + − + ...
2 3
so we consider the integral Z
(z−1)2 3
− (z−1)
e(1+ 2 3 +...)
dz.
C ′′
dz
We parameterise the path of integration as z = 1 + iv, −δ ≤ v ≤ δ and use the chain rule to write dz = dv dv = ivdv and the
integral becomes: √
Z v=δ Z v=δ Z v=∞
2
γ(1− v2 − (iv)
3
) γ(1− v2
) γ −γv2 ieγ 2π
i e 3 dv ∼ i e 2 dv ∼ ie e 2 dv = √
v=−δ v=−δ v=−∞ γ
where we have used Laplace’s method and the usual trick involving the integration path. So the leading order estimate is:
1 eγ 1
∼ √ γ 2 −γ .
Γ(γ) 2π
More terms To deform the path of integration into the steepest descent path, we make the exact change of variables
f (z) = 1 − η 2 .
Again, this change of variables is motivated by the behaviour of the function f (z) = z − ln z. We want a curve, parameterized
by η, passing through z = 1 and η = 0, with ℑ(f ) = 0, ℜ(f ) < 1 and with locally quadratic behaviour in η. With this change
of variables the integral becomes
Z
1 1 γ 1−γ ∞ −γη2 dz
= e γ e dη
Γ(γ) 2πi −∞ dη
dz
which is still exact. However, to find dη we must make some approximations. We first invert the change of variables to find
dz
z = z(η) and then differentiate to find dη . We determine that
√ 2 i 2 4 i 4 6
z(η) = 1 + i 2η − η 2 − √ η 3 − η + √ η5 − η + ··· .
3 9 2 135 540 2 8505
152
dz
Now that z(η) is known, we also know dη in power series form, and although it has only a finite radius of convergence, we
invoke Watson’s Lemma to integrate term by term. The result is
Z
1 1 γ 1−γ ∞ −γη2 η2 η4
= √ e γ e 1− + + · · · dη
Γ(γ) π 2 −∞ 6 216
γ 1/2
γ −γ 1 1
∼e γ 1− + + ··· .
2π 12γ 288γ 2
This answer is in agreement with Stirling’s formula found earlier.
153
1
0.8
0.6
0.4
0.2
0 1 2 3 4 5
Figure 62: The error function erf(x)
10 Appendix
10.1 Gamma Function
Z ∞
Γ(x) ≡ ux−1 e−u du, (x > 0) (747)
0
√
Note that Γ(x + 1) = xΓ(x), Γ(1/2) = π and Γ(0) = ∞.
10.2 Error function and complementary error function

Z x Z ∞
2 −s2 2 2
erf(x) ≡ √ e ds; erf c(x) ≡ 1 − erf(x) ≡ √ e−s ds (748)
π 0 π x
erf(−x) = erf(x); erf(0) = 0; erf(∞) = 1; erf c(0) = 1; erf c(∞) = 0. (749)
2
e−x
erf(x) ∼ 1 − √ (x → ∞) (750)
πx
10.3 Airy Functions

See Abramowitz and Stegun (1972) for more extensive results.
10.3.1 Airy Differential equation
′′
y = xy, for − ∞ < x < ∞. (751)
10.3.2 General solution
y(x) = αAi(x) + βBi(x) (752)

where
∞
X
1 Γ( k+1 ) sin( 2π (k + 1))
Ai(x) ≡ 3 3
(31/3 x)k (753)
32/3 π k=0
k!
1 ′ 1
= Ai(0)(1 + x3 + . . .) + Ai (0)(x + x4 + . . .) (754)
6 12
154
8
0.5 1 1.5 2 2.5 3 3.5 4
Figure 63: Gamma function for 0 < x ≤ 4.
2.0
Ai(x)
1.5
Bi(x)
Airy Functions
1.0
0.5
0
-0.5
-10 -8 -6 -4 -2 0 2
x - axis
Figure 64: Airy functions.
155
and
Bi(x) = eπi/6 Ai(xe2πi/3 ) + e−πi/6 Ai(xe−2πi/3 ) (755)

1 ′ 1
= Bi(0)(1 + x3 + . . .) + Bi (0)(x + x4 + . . .). (756)
6 12
10.3.3 Particular values
Γ(1/3) ′ 31/6 Γ(2/3) √ ′ √ ′

Ai(0) = 1/6
, Ai (0) = − , Bi(0) = 3Ai(0), Bi (0) = − 3Ai (0)
2π3 2π
Z ∞ Z 0 Z 0
1 2 2π
Ai(x)dx = , Ai(x)dx = , Bi(x)dx = 0, Γ(1/3)Γ(2/3) = √ .
0 3 −∞ 3 −∞ 3
10.3.4 Asymptotic expressions for large

x
For ζ = 32 |x|3/2 ; η = 5
72ζ ; ν = 7
72ζ
1 π π
Ai(x) ∼ √ 1/4
[cos(ζ − ) + η(x) sin(ζ − )]
(x → −∞) (757)
π |x| 4 4

1 1 2
∼ √ 1/4 exp (−ζ)[1 − η(x)] ∼ √ 1/4 exp − x3/2 (x → ∞) (758)
2 πx 2 πx 3
1 π π
Bi(x) ∼ √ 1/4
[cos(ζ +
) + η(x) sin(ζ + )] (x → −∞) (759)
π |x| 4 4

1 1 2 3/2
∼ √ 1/4 exp(ζ)[1 + η(x)] ∼ exp x (x → ∞) (760)
πx πx1/4 3
1/4
′ 1 |x| π π
Ai (x) ∼ √ [cos(ζ + ) − ν(x) sin(ζ + )] (x → −∞) (761)
π 4 4
−1x1/4
∼ √ exp(−ζ)[1 + ν(x)] (x → ∞) (762)
2 π
1/4
′ |x| π π
Bi (x) ∼ √ [cos(ζ − ) − ν(x) sin(ζ − )] (x → −∞) (763)
π 4 4
x1/4
∼ √ exp(ζ)[1 − ν(x)] (x → ∞) (764)
π
10.3.5 Useful Taylor (Maclaurin) series

When using asymptotic methods, these are especially useful when x ≪ 1, when the truncated series can be very accurate.
Note the use of the order symbols below.
156
x2 ′′
′
f (0 + x) = f (0) + xf (0) + f (0) + O(x3 )
2!
exp(x) = 1 + x + x2 /2! + x3 /3! + O(x4 ) (765)
2 3 4 5
ln(1 + x) = x − x /2 + x /3 − x /4 + O(x )(−1 < x ≤ 1) (766)
3 5 7
sin(x) = x − x /3! + x /5! + O(x ) (767)
2 4 6
cos(x) = 1 − x /2! + x /4! + O(x ) (768)
3 5 7
tan(x) = x + x /3 + 2x /15 + O(x )(| x |< π/2) (769)
3 5 7
sinh(x) = x + x /3! + x /5! + O(x ) (770)
2 4 6
cosh(x) = 1 + x /2! + x /4! + O(x ) (771)
3 5 7
tanh(x) = x − x /3 + 2x /15 + O(x )(| x |< π/2) (772)
3 5 7
2 x x x
erf(x) = √ (x − + − + . . .) (773)
π 3.1! 5.2! 7.3!
10.4 Binomial series

n
X
n n n−k k n n−1 n(n − 1) n−2 2 n n 0 n n−1 2
(a + x) = a x = a + na x+ a x + ... = a x + a x + ... (774)
k 2! 0 1
k=0

n n(n − 1) 2 n 0 n 2
(1 + x) = 1 + nx + x + ... = x + x + . . . ( converges for |x| < 1 for non integer n.) (775)
2! 0 1
The last two results can be written concisely as

∞
X Γ(m + 1)xm
(1 + x)n = (776)
m=0
Γ(n + 1 − m)m!
When n is a positive integer, the series truncates (after m = n all terms are zero) by virtue of the fact that Γ(n) → (−1)n ∞
as n → 0, −1, −2, −3 . . .
10.5 Useful Inequalities

Z Z
b b

f (x)dx ≤ |f (x)| dx ≤ (b − a)maxx∈[a,b] |f (x)| .
a a
This is best illustrated by drawing some arbitrary function f (x) taking on both positive and negative values on [a, b]. For
example let f (x) = sin x on [0, 2π] and examine this result graphically (see figure 10.5).
10.5.1 Cauchy Schwartz

For two (possibly n-dimensional) vectors this is easy
|a.b|2 ≤ |a|2 |b|2 because a.b = |a| |b| cos θ and |cos θ| ≤ 1.
For functions, if we define an inner product via:
Z b
(f, g) ≡ f (x)g(x)dx
a
then Cauchy-Schwartz is similar to the vector case i.e.,
Z !2 Z Z
b b b
f (x)g(x)dx ≤ f (x)2 dx g(x)2 dx.
a a a
157
Figure 65: Clearly the area under y = | sin x| is less than that under y = 1.
10.5.2 Fredholm alternative

Fredholm alternative theory for matrices For the matrix equation Ax = b where A is m × n (not necessarily square).
1. The solution of Ax = b (if it exists) is unique if and only if the only solution of Ax = 0 is x = 0.
2. The equation Ax = b has a solution if and only if (b, v) = 0 for every vector satisfying A∗ v = 0.
The latter of these two statements is called the Fredholm alternative theorem. Here A∗ = ĀT is the adjoint matrix where
the bar denotes the complex conjugate. For a real matrix A∗ = AT .
Fredholm alternative theorem for differential operators (see Stakgold) For the nth order boundary value problem
Lu = f, a < x < b with associated boundary conditions B1 u = γ1 . . . Bn u = γn and the associated adjoint problem L∗ v = 0
with homogeneous boundary conditions B1∗ v = 0 . . . Bn∗ v = 0,
1. the solution of Lu = f is unique if and only if Lu = 0 has only the trivial solution u = 0.
2. The solution of Lu = f exists if and only if (f, v) = 0 for all v in the null-space of L∗ i.e., for all v satisfying the
associated homogeneous adjoint problem.
It follows that Lu = f can be solved for any fˆ if and only if L∗ v = 0 =⇒ v = 0.
158
Perturbation methods Sheet 1
Assume that ε ≪ 1 and all other constants are O(1).
1. Determine 3 terms in the asymptotic expansions for:

√ p
(a) (1 − 38 a2 ε + 256
51 4 2 −1
aε) (b)cos( 1 − εt) (c) 1 − 2ε + 2ε2 (d) sin(1 + ε − ε2 ). (N1.1)
p
2. Let µ = µ0 + ε µ1 + ε 2 µ2 in h = 23 (1 − 1 − 3µ(1 − µ), expand for small ε and keep 3
terms. (N1.3)
3. Determine the order (O and ∼) of the following functions: (a) sinh(1/ε ) (b) ln(1 + sin ε)
(c) ln(2 + sin ε ) (d) exp(ln(1 − ε)). (N1.4)
4. (BO79) Which of the following statements is true? (Show your reasoning).
x ≪ 1/x(x → ∞), x1/2 ≪ x1/3 (x → 0), ln x ≪ x(x → ∞), (1)

5 1/4 x x 2
(ln x) ≪ x (x → ∞), e + x ∼ e (x → ∞), x ∼ x(x → 0). (2)
5. (N1.8) Arrange in descending order (using ≫) so that it forms a well ordered set of gauge
functions:
√ 1 3
ε2 , ε, ln(ln(1/ε)), 1, ε 2 ln(1/ε), ε ln(1/ε), exp(−1/ε), ln(1/ε), ε 2 , ε, ε2 ln(1/ε) (3)
6. Give an example to show that f = O(g) as ε ↓ ε 0 (and ε → ∞) does not necessarily

mean that exp(f ) = O(exp(g)). (H6/2).
7. Verify that if f (x) ∼ g(x) and g(x) is bounded as x → x0 , then ef ∼ eg .
8. Prove that exp(−1/ε) ≪ εα , for any α. (It may be useful to put ε = 1/δ and consider the
limit δ → ∞ while using l’Hospital’s rule).
9. Let f (x, ε) ≡ f1 +f2 +f3 . If φ1 ≡ f1 , φ2 ≡ f1 +f2 , φ3 ≡ f1 +f2 +f3 show that φ1 ∼ f, φ2 ∼ f

and φ3 ∼ f uniformly for all x in some interval I means that f1 ≪ f2 ≪ f3 uniformly.
1 √
10. If α + βε + γε2 ∼ (2 + ε) 2 , show that α = 2, β = 2√1 2 , γ = − 161√2 .
11. Prove l’Hospital’s rule for two functions f (x), g(x) which are analytic (i.e., which have
Taylor series which converge to the function) in the case where f (0) = 0, g(0) = 0 and we
wish to evaluate limx→0 fg(x)
(x)
.
12. Find an asymptotic expansion for f (ε) = sin eε . (Do this in two steps: first expand eε .
Then, in effect, expand sin (1 + y), y ≪ 1).
1
Perturbation methods Sheet 2 (Algebraic equations; regularly perturbed ODEs)
1. Determine 2 term leading order expressions for each root of: x2 + ε x − 1 = 0 and
εx2 + 2x − 1 = 0. What is the significant difference between these two perturbation
problems? Check your answers by finding the exact solution and expanding for small ε.
2. Find a two-term asymptotic expansion, for small ε, of the following:

(a) The point xm = (xm , ym ) at which the function f (x, y) = x2 + 2ε sin(x + ey ) + y 2
2
attains its minimum value. (Recall the conditions for a minimum: fxx fyy − fxy > 0 and
fxx > 0 and fyy > 0).
(b) The perimeter of the two-dimensional golf ball described as r = 1 + ε cos(20θ), where
0 ≤ θ ≤ 2π. Explain why ε needs to be fairly small before this two-term expansion can
be expected to produce an accurate approximation of the arc length. Also, in geometric
terms, explain why this happens. (Hint: ds2 = dr 2 + r 2 dθ2 ).(H24/3)
3. (N2.1)Determine 2 terms in the expansion of each root of
x3 − (2 + ε)x2 − (1 − ε)x + 2 + 3ε = 0, x3 − (3 + ε)x − 2 + ε = 0.

√ p
(Answer: 1 + 3ε/2, −1 − ε/6, 2 − ε/3). (Answer: 2 + ε/9, −1 ± ε 2/3).
4. (N2.2) Determine 2 terms in the expansion of each root of:
ε(u3 + u2 ) + 4u2 − 3u − 1 = 0; εu3 + u − 2 = 0.

3
(Answer: u = −1/4 + ε 320 , 1 − ε 25 , −4/ε − 7/4). (Answer: u = 2 − 8ε, ± √iε − 1).
π
5. Find the root of ǫx tan x = 1 which is close to π/2. (Answer: x ∼ 2
− ε π2 .)
6. Determine 2 term expansions for the large roots of x tan x = 1; x cot x = 1.(N2.4)
7. Show graphically that x exp(−x) = ε has a small root and a large one. Hence show that
the two approximate roots are
x1 ∼ ε + ε2 , x2 ∼ ln (1/ε) + ln (ln (1/ε)).
Hint: Remember that if ε ≪ 1, then ln (1/ε) ≡ − ln ε ≫ 1. Check this numerically with

Maple. It is necessary to take very small values of ε to get good agreement because of
the ln ln term.
8. A well-studied problem in solid mechanics concerns the deformation of an elastic body

when compressed by a rigid punch. This gives rise to having to evaluate the following
integral (Gladwell, 1980): Z ∞
I= N(x) sin(λx)dx,
0
where
2α sinh(x) − 2x
N(x) = −1
1 + α2 + x2 + 2α cosh(x)
and α, λ are constants with 1 ≤ α ≤ 3. The indefinite interval complicates finding a
numerical value of I, and so, write
Z x0
I= N(x) sin(λx)dx + R(x0 ).
0
2
The objective is to take x0 large enough that the error R(x0 ) is relatively small.
(a) Find a two-term expansion of N(x) for large x.
(b) Use the first term in the expansion you found in part (a) to determine a value of x0
so |R(x0 )| ≤ 106 .
(c) Based on the second term in the expansion you found in part (a), does it appear that
the value of x0 you found in part (b) is a reasonable choice? (H18/10)
9. (N1.14) In terms of x, where does each of the following series fail to be uniformly well
ordered asymptotic series and what are its regions of non-uniformity?
3
(a) x 2 + ε (ln 1+x
1−x
) + ... with x ∈ [0, 1)..
1+x 2
(b) ε cos x + ε 2 1−x 2 cos 2x + ... with x ∈ [0, 1).
√ 2
(c) x2 − 1 − √εx x2 −1
+ ... with x ∈ [0, 1).
(d) 1 − ε x + ε 2 x2 + ... ∀x.
(e) sin x + ε cos x + ... with x ∈ [0, 2π].
Now consider just the first two terms of each of the above expressions and regard each
of these as a function f (x, ε) and not as an asymptotic series. Is the first term (labelled
g(x, ε) ) a uniformly valid approximation for the first two terms i.e., is limε→0 fg = 1
3
uniformly for all x? If not where does it fail? (For example in (a) is g(x) = x 2 a
3 1+x
uniformly valid approximation to f (x, ε) = x 2 + ε (ln 1−x )?)
10. In the projectile problem in the notes, to account for air resistance one obtains the
equation
d2 x gR2 k dx
2
= − 2
− ,
dt (x + R) R + x dt
where k is a nonnegative constant. Assume here that x(0) = 0 and x′ (0) = v0 .
(a) What is the nondimensional version of this problem if one uses the same scaling as in
the notes?
(b) Find a two-term asymptotic (regular perturbation) expansion of the solution for small
ε. Assume in doing this that α = kv0 /gR is independent of ε.
(c) Does the addition of air resistance increase or decrease the flight time?(H31/3)
11. Find a regular 2 term asymptotic expansion for the solution of:
′′ ′
(a) y +εy − y = 1; y(0) = y(1) = 1.
′′
(b) y +y + εy 3 = 1; y(0)0, = y(π/2) = 1.
′′
(c) y +f (εy) = 0; y(0) = y(1) = 0, (f (s) is a smooth function). (H31/1)
12. Find two term expansions for each of the solutions of εx3 − 3x + 1 = 0.
13. Graph the function x sin x on [−3π, 3π] and verify that
√ the function is even. Hence
ε
deduce the approximate roots of x sin x = ε. (Answer:± ε, ±(Nπ − N π ), N odd integer,
± (Nπ + Nεπ ), N even integer).
3
Perturbation methods Sheet 3 (Matched asymptotic expansions)
Assume that ε ≪ 1 and other constants are O(1).
1. Consider the system of ODEs:

du dv
= (u − v) + u; ε = u − v.
dx dx
If the initial conditions are u(0) = 1, v(0) = 1, show that u(x) ∼ v(x) ∼ ex . Would you
expect a boundary layer if v(0) = 0? Solve the problem in this case by using matched
asymptotic expansions to leading order. Check your answer by finding the exact solution
using Maple.
2. Use matched asymptotics to find a leading order solution to the boundary value problem:
εy ′′ + y ′ = 0; y(x = 0) = 1; y(x = 1) = 2.
Sketch the approximate solution. Verify your answer by solving the equation exactly and
taking appropriate limits as ε → 0.
3. A boundary layer at infinity. Consider the problem:
u′′ + εu′ = 0, u(0) = 0, u(∞) = 1.
Show that the leading order straightforward solution is of the form u0 = Cx for some
constant C. Rescale at large distances via x = x̄/ε anticipating a boundary layer at
infinity. Show that matching requires that C = 0. Examine the exact solution and
reconcile it with the matched asymptotics approach. (See also Viscous Flow, Ockendon
and Ockendon, p.51).
4. Consider the boundary value problems:

′′ ′
εy + y + y = 0; y(0) = α, y(1) = β. (4)
′′ ′
εy − y + y = 0; y(0) = α, y(1) = β. (5)
′′ ′
εy − y = 1; y(0) = α, y(1) = β. (6)
Determine the exact solution.Use the method of matched expansions to determine a first
order uniform expansion. Compare your answer with the exact solution. (N12.1,N12.2,N12.3)
5. A small parameter multiplying the highest derivative does not guarantee that boundary
or interior layers are present. This problem presents examples of such situations.
(a) After solving each of the following problems, explain why the method of matched
asymptotic expansions cannot be used (in a straightforward manner) to find an
asymptotic approximation of the solution.
i. ε2 y ′′ + ω 2 y = 0, for 0 < x < 1 and ω 6= 0.
ii. ε2 y ′′ + εy ′ + y = 1, for 0 < x < π, where y(0) = y(π) = 0.
iii. ε2 y ′′ = y ′ , for 0 < x < 1, where y ′ (0) = −1 and y(1) = 0.
iv. ε2 y ′′ − 2xy ′ − 2y = 0, for −1 < x < 1, where y(−1) = 1 and y(1) = 2.
(b) For what values of α and β, if any, will the following problem have a boundary
layer(s)?
ε2 y ′′′′ + ω 2 y = 0, for 0 < x < 1,
where y(0) = y(1) = 0, y ′′(0) = α, y ′′(1) = β. Also, ω is a positive constant. (H 58/3)
4
6. The Michaelis-Menten reaction scheme for an enzyme catalyzed reaction is
ds dc
= −s + (s + k − 1) c and ε = s − (s + k) c, for t > 0,
dt dt
where s(0) = 1 and c(0) = 0. Here s(t) is the concentration of substrate, c(t) is the
concentration of the chemical produced by the catalyzed reaction, and k is a positive
constant. Find the first term in the expansions in (i) the outer layer, (ii) the initial layer
e
t = t/ε , and (iii) a composite expansion.(H 59/10)
7. Consider the problem

′′′ ′ ′
εy − y = 1; y(0) = α, y (0) = β, y(1) = γ. (7)
Determine the exact solution and use it to show there is a boundary layer at each end.
Determine a two term uniform expansion and compare √ your answer with the exact solu-
tion. (Hint: look for expansions of the form y = y0 + εy1 in each region). (N12.12)

′′′ ′ ′
εy − (2x + 1)y = 1; y(0) = α, y (0) = β, y(1) = γ. (8)
Determine a 2 term uniform expansion

√ (with boundary layers at each end). (Hint: look
for expansions of the form y = y0 + εy1 in each region). (N12.14)

′′ ′
εy − yy − xy = 0; y(0) = 1, y(1) = 1. (9)
Determine a first order uniform expansion for the case y = O(1) (so there is no need to
scale y) with an interior layer in [0, 1].(N12.21)
10. A ‘small’ mass m hangs from a weightless spring with internal damping proportional to
the speed. A known vertical impulse I (instantaneous momentum change) is imparted to
the mass by striking it with a hammer. Initial conditions on the vertical deflection y ∗ (t∗ )
can thus be taken to be:
dy ∗
y ∗ (t∗ = 0) = 0; m ∗ (0) = I.
dt
The governing differential equation for the motion of the mass is:
d2 y ∗ dy ∗
m + µ + ky ∗ = 0,
dt∗2 dt∗
where µ, k are the damping and string constants respectively. (What are the respective
units?) We assume that the mass in sufficiently small that we have a strongly overdamped
(i.e., the spring is a ‘strong’ one) situation and the mass will quickly return to rest after the
impulse is expended in stretching the spring. Show that a certain choice of dimensionless
variables reduces the problem to:
d2 y dy dy
ε + + y = 0; y(0) = 0, ε (0) = 1, ε ≪ 1.
dt2 dt dt
11. Solve, using matched asymptotics, noting that there are boundary llayers at each end of
the interval:
ε2 y ′′ + εxy ′ − y = −ex , for 0 < x < 1,
5
where
y(0) = 2 and y(1) = 1.
Note that there is a similar problem in the notes, but the term on the right hand side is
different.
12. Solve the singular perturbation problem (write down an associated inner and outer prob-
lem)
d2 y dy dy
ε 2 + + y = 0; y(0) = 0, ε (0) = 1.
dt dt dt
Do not impose any initial condition on the outer problem. Then find an inner solution
(t = εT ) which satisfies both initial conditions. Complete the determination of the outer
solution by matching the two solutions. Check your answer by finding the exact solution
of this linear problem. Can you explain why the original scaling represents the true
balance in the problem except for small times?
13. (H82/1)Find a composite expansion (with corner layer) for the following problem:
′′ ′
εy + y(y )2 − y = 0; y(0) = 1, y(1) = 1. (10)
(For this problem go to O(ε) in both the inner and outer solutions. Check your answers
by using Maple to generate and plot a numerical solution. You will need the following
commands:
restart; epsilon:=0.1;ode:=epsilon*diff(y(x),x,x)+y(x)*diff(y(x),x)ˆ2-y(x);
bcs:=y(0)=1,y(1)=1;t1:=dsolve({ode,bcs},numeric};plots[odeplot](t1);)
14. Find a composite expansion (with corner layer) for the following problem:
′′ ′
εy = 9 − (y )2 ; y(0) = 0, y(1) = 1. (11)
15. A simple toy model for spin-coating of a liquid on a rotating substrate with evaporation
effects included is:
dL∗
= −αL∗2 − β; L∗ (0) = L0 (12)
dt∗
2
where L∗ (t∗ ) is the film thickness, t∗ is the time, α(= 2ω
3ν
) is a parameter representing
the relative importance of spin effects and viscosity and β is an evaporation parameter.
Though the initial liquid distribution will usually be spatially dependent, the film thick-
ness quickly becomes spatially uniform and can thus be considered to be only a function
of the time. Motivated by the physical spin-coating problem, typical parameter values
are:
α ∼ 108 m−1 s−1 ; β ∼ 10−6 ms−1 ; L0 ∼ 10−6 m (13)
The terms on the right hand side of (12) can be thought of as representing the tendency
of the liquid film to thin out via two processes: spinning/viscous (flow) effects (as repre-
sented by α) and evaporation (as represented by β) respectively. Initially when the film
is relatively thick the film thickness L∗ (t∗ ) decreases predominantly via flow. As the film
thickness becomes sufficiently small, viscous effects become relatively more important,
the flow slows down, and eventually the effects of evaporation take over. First scale the
problem via
6
θ θ
21
∗ β 1
L =l ; t∗ = t 1
α (αβ) 2
to get
dl 1
= −l2 − 1; l(0) = 1
dt ε2
β
with ε ≡ αL20
≪ 1. Now consider the rescaling
1 1
l= 1 L(T ); t = ε 2 T
ε 2
and obtain a leading order solution using matched asymptotic expansions.

16. The shape of liquid meniscus near a solid vertical glass cylinder placed in an infinite sea
of water
is modelled by: !
d2 z dz
dr 2 dr
γ 3 + 1 − ρgz = C
(1 + ( dz
dr
)2 ) 2 r(1 + ( dz
dr
)2 ) 2
where z = z(r) is the meniscus height above the undisturbed flat liquid level, r is distance
from the axis of symmetry, γ(= 0.08 Nm−1 for water) is the surface tension, ρ ( 103 Kg
m−3 for water) is the liquid density, g(ms−2 ) and C is a dimensional constant which we
can set to zero by setting the plane z = 0 to coincide with the undisturbed meniscus
height at ∞. The boundary conditions on the cylinder is:
dz
= tan θ at r = R,
dr
where θ is the contact angle and R (m) is the cylinder radius. At infinity we require that
z and all its derivatives tend to zero.
q
γ
Note that there are two length scales in this problem: R and ρg . The latter is referred
to as the capillary length: for water and it is of the order of a couple of mm (check this).
One can non-dimensionalise this problem by scaling lengths with R: (z = Ry, r = Rx)
which leads to:
d2 y dy
dy ′
dx2
dy 2 23
+ dx
dy 2 12
− ǫ2 y = 0, (x = 1) = tan θ, y, y (x) → 0 as x → ∞
(1 + ( dx )) x(1 + ( dx )) dx
2
Treat this as the inner scaling of a boundary layer problem where ǫ2 = ρgR
γ
≪ 1, i.e.,
we are dealing with cylinders whose radius R is much smaller than the capillary length
7
q
γ
ρg
(Note that to solve the leading order problem one must solve a Bernoulli ODE). The
outer problem is obtained via the rescaling (x = X/ǫ, y(x) = Y (X)) which leads to:
d2 Y dY
dX 2 dX dY tan θ ′
dY 2 32
+ dY 2 21
−Y = 0, (X = ǫ) = , Y, Y (X) → 0 as X → ∞.
(1 + ǫ2 ( dX ) ) X(1 + ǫ2 ( dX )) dX ǫ
(The leading order outer problem is the modified Bessel equation of zero order). Now
solve the problem using matched asymptotic expansions (to leading order).
17. Solve the advection-diffusion problem
ut + ux = εuxx ; u(0, t) = 1; u(1, t) = 2, u(x, 0) = 0; ε ≪ 1,
by exploiting the smallness of ε to solve an advection problem on most of the domain.

Show in detail how to resolve the boundary layer near x = 1. The advective part of the
solution should be in the form of a moving step function. Show how this too can be
smoothed by insertion of a boundary layer.
8
Perturbation methods Sheet 4 (Multiple scales)
1. (M238)Show that a regular perturbation approach fails for the initial value problem:
′′ ′ ′
y + 2εy + (1 + ε)y = 0; y(0) = α, y (0) = 0. (14)
Find a leading order uniformly valid solution using multiple scales.

2. (N12.1-12.3) Solve the boundary layer problems in question 1 of the previous sheet using
multiple scales.
3. (H113/1) Find a first term expansion of the solutions of the following which is valid for
large t:
(a)
′′ ′ ′
y + ε(y )3 + y = 0, y(0) = 0, y (0) = 1, (t > 0). (15)
(b)
′′ ′ ′
εy + εy + y = cos t, y(0) = 0, y (0) = 0, (t > 0). (16)
4. Consider the van der Pol oscillator problem (with small damping):
′′ ′ ′
y − ε(1 − y 2)y + y = 0, y(0) = 1, y (0) = 0, (t > 0).
Using multiple scales show that the leading order solution is
2
y ∼ A(εt) cos t; A(εt) = 1 .
(1 + 3 exp(−εt)) 2
What happens to the amplitude of the oscillation for large times?

5. Consider Rayleigh’s equation:
1
ÿ + ε( ẏ 3 − ẏ) + y = 0.
3
(a) Obtain a regular perturbation solution and show that this contains secular terms.
(b) Using multiple scales (t1 = t, t2 = εt) show that the leading order solution is
1
y = 2 cos(t + β)(1 + ke−εt )− 2 ,
where k, β are constants of integration.
Why is this solution an improvement on the regular perturbation solution?
You may use the result
dA 1 2
2 + A(τ )3 − A(τ ) = 0 has solution A = √ , k constant.
dτ 4 1 + ke−τ
6. Noting that the initial value problem
′′ ′ ′
εy + y + y = 0, y(0) = 0, y (0) = 1, (t > 0). (17)
has an initial (boundary layer) of O(ε), solve the problem to leading order using matched
asymptotics. It will be necessary to rescale the dependent variable (y = εY ) to do so
(because of the inhomogeneous initial condition). Now use multiple scales to solve the
same problem. (Hint: what are the important time scales as deduced by using matched
asymptotics?)
9
Perturbation methods Sheet 5 (WKB method)
1. Consider the eigenvalue problem:

1
y ′′ + f (x)y = 0; y(1) = 0; y(2) = 0
ε2
1
where the the eigenvalues are ε2
.
(a) Solve the boundary value problem if f (x) ≡ 1 and write down the eigenvalue equation
(without solving it).
(b) If f (x) = x2 , use a WKB approximation to find a leading order approximation for
y(x) for ε ≪ 1.
1
(c) Show that a leading order approximation for the large eigenvalues is given by: εn
=
2
3
nπ, n ∈ Z + .
2. Show that the large eigenvalues of

′′
u + λ2 f (x)u = 0, (f (x) > 0); u(0) = 0, u(1) = 0 (18)
hR p i−1
1
are given by: λn = nπ 0
f (x) . (N14.6)
3. Show that the large eigenvalues of

′′ ′
u + λ2 f (x)u = 0, (f (x) > 0); u(0) = 0, u (1) = 0 (19)
hR p i−1
1 1
are given by: λn = (n + 2
)π 0
f (x) . (N14.6)
4. Find an approximate solution of the boundary value problem:
′′
εy + y = 0; y(0) = 0, y(1) = 1 (20)
for ε ≪ 1, using the WKB approach. (BO487/2)
5. By consideration of the problem

d2 y
ε2
= xy,
dx2
find an asymptotic solution of the Airy equation (y ′′ = xy) for large positive x.
(To solve the Airy equation (y ′′ = xy) for large x carefully, one can use an artificial small
2 1
parameter ε and scale x ∼ ε− 3 and y ∼ ε 6 to transform the equation to the form at the
1 2 3
beginning of the question. Independent solutions for large x are x− 4 e± 3 x 2 ).
6. (Nayfeh, p.386) Apply the WKB method to the problem:

1
xy ′′ + y ′ + x(1 − x2 )y = 0.
ε2
10
Perturbation methods Sheet 6 (Integrals)
1. (N3.3) Use integration by parts to show that

Z ∞
exp(−t) 1 1! 2! 3!
dt ∼ exp(−x)( − 2 + 3 − 4 )(x ≫ 1) (21)
x t x x x x
2. (L94) Consider the integral

Z ∞
sin t −λt
I(λ) = e dt, λ ≫ 1. (22)
0 t
Using Maple, graph the integrand for increasing values of λ(= 1, 5, 10..). Use Laplace’s
method to get an asymptotic approximation to the integral and show that
1 2! 1
I(λ) ∼ − 3
+ O( 5 ). (23)
λ 3!λ λ
Use Maple to obtain a numerical approximation and compare this with the asymptotic
estimates. You will have to use an approximation for the upper limit in the integral.
3. Use integration by parts to show that

Z λ 2
2 eλ
es ds ∼ , (λ ≫ 1).
0 2λ
Verify this result by using Laplace’s method. Use Maple (or otherwise) to test the result
by evaluating the integral numerically.
4. (L100) Show that the leading order asymptotic behaviour of

Z π
I(λ) = eλ sin t dt, λ ≫ 1 (24)
0
q
2π
is I(λ) ∼ eλ λ
. Deduce a leading order estimate for:
Z 2nπ
I(λ) = eλ sin t dt, λ ≫ 1, n ∈ Z + . (25)
0
5. For x ≫ 1, use Laplace’s method to show that

Z 2π √
2 x sin2 t 3 πex
cos te dt ∼ .
π
2
4x3/2
11

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Module MS4407: Perturbation techniques and asymptotics

Prof. Stephen O’Brien.

Objectives and Assessment

3 Algebraic and Transcendental Equations 20

4 Regular Perturbation Solution of Differential Equations 28

5 Matched Asymptotic Expansions 42

7 Method of strained co-ordinates 106

8 The WKB Method 110

9 Asymptotic expansion of integrals 132

In section 4.6 we review solution techniques for some simple ODEs.

1.1.1 The wave equation

Figure 1: An oscillating elastic string

x = lx∗ , u = u0 u∗ , t = (l/c)t∗ , (5)

∂u ∂(u0 u∗ ) u0 ∂u∗ ∂2u u ∂ 2 u∗ ∂u ∂(u0 u∗ ) u0 c ∂u∗

f (x) ∂u∗ (x∗ , 0) g(x)l

1.1.2 Scaling and rescaling

1.2 Example: projectile problem

1.2.1 Scaling the problem

2.2 Useful results

2.2.1 Theorem 1.1 (Taylor series with remainder)

2.2.2 Theorem 1.2 (l’Hospital’s rule)

2.3 Order Symbols

1. f = O(φ) as ε ↓ ε0 if the limit

|f (ε)| ≤ k0 |φ(ε)| for ε0 < ε < ε1 .

|f (ε)| ≤ δ|φ(ε)| for ε0 < ε < ε2 .

This is also commonly written as f ≪ φ.

4. If g 6= 0, then f = O(g) as ε → ε0 if and only if f /g is bounded in a (punctured) neighborhood of ε0 , and f = o(g) if

2.4 Asymptotic Approximations

2.5 Asymptotic Expansions (Series) and Gauge Functions

f = a1 φ1 (ε) + o(φ1 ) and f = a1 φ1 (ε) + a2 φ2 (ε) + o(φ2 ) etc.

In terms of limits this means that

1. φ1 = εα , φ2 = εβ , φ3 = εγ , · · · where α < β < γ < · · · (for ε ↓ 0), and

Asymptotic expansion of a known function f (ε)

Example using Taylor’s theorem

0.0 0.1 0.2 0.3 0.4 0.5 0.6 0.7

Example using l’Hospital’s rule

where the an are independent of ε. This series converges if and only if

RN (ε) = o(εN ) (ε → 0), (37)

1·32 ε2 1·32 ·52 ·72 ε4

2.6 Stokes phenomenon in the complex plane

as z → ∞ with π4 < arg z < 3π 5π 7π

2.6.1 Manipulating Asymptotic Expansions

f (x, ε) ∼ a1 (x)φ1 (x, ε) + a2 (x)φ2 (x, ε) as ε ↓ ε0 . (42)

f ∼ 0 + 0·ε + 0·ε2 . . . . (44)

f (x, ε) ∼ a1 (x)φ1 (ε) + a2 (x)φ2 (ε) as ε ↓ ε0 , (45)

3 Algebraic and Transcendental Equations

(1 + εx1 )5 − ε(1 + εx1 ) − 1 = 0 (51)

1 + 5εx1 − ε − 1 + O(ε2 ) = 0 (52)

x20 + 2εα x0 x1 + . . . + 2ε(x0 + εα x1 + . . .) − 1 = 0. (57)

3.3 Example (A singular perturbation)

3.3.1 Rescaling the equation

ε1+2ν X 2 + 2εν X − 1 = 0. (61)

(X02 + 2εα X0 X1 + . . .) + 2(X0 + εα X1 + . . .) − ε = 0. (64)

ε0 X02 + 2X0 = 0 =⇒ X0 = 0, −2 and α = 1

x20 + 2x0 x1 εα + . . . + 1 + εx0 + . . . = 5.

Figure 8: Graphs of y = x2 + exp(εx) and y = 5 for ε = 0.1.

ε(π/2 + εx1 ) tan (π/2 + εx1 ) ∼ 1 + O(ε2 ).

Figure 9: The exponential functions.

µ + εsech(−ε−1 + µ/ε) = 0. (72)

Figure 10: Graph of y =sech x.

Figure 11: The graph of y = x + 1 + εsech ( xε ).

x ∼ −1 − 2εe−1/ε + ν(ε), (74)

3.7 Iterative methods